license: creativeml-openrail-m
base_model: stabilityai/stable-diffusion-xl-base-1.0
tags:
- sdxl
- sdxl-diffusers
- text-to-image
- image-to-image
- diffusers
- simpletuner
- not-for-all-audiences
- lora
- controlnet
- template:sd-lora
- standard
pipeline_tag: text-to-image
inference: true
widget:
- text: unconditional (blank prompt)
parameters:
negative_prompt: blurry, cropped, ugly
output:
url: ./assets/image_0_0.png
- text: A photo-realistic image of a cat
parameters:
negative_prompt: blurry, cropped, ugly
output:
url: ./assets/image_1_0.png
- text: prompt not found (2)
parameters:
negative_prompt: blurry, cropped, ugly
output:
url: ./assets/image_2_0.png
- text: prompt not found (3)
parameters:
negative_prompt: blurry, cropped, ugly
output:
url: ./assets/image_3_0.png
simpletuner-controlnet-sdxl-lora-test
This is a ControlNet PEFT LoHa derived from stabilityai/stable-diffusion-xl-base-1.0.
The main validation prompt used during training was:
A photo-realistic image of a cat
Validation settings
- CFG:
4.2
- CFG Rescale:
0.0
- Steps:
20
- Sampler:
ddim
- Seed:
42
- Resolution:
1024x1024
Note: The validation settings are not necessarily the same as the training settings.
You can find some example images in the following gallery:

- Prompt
- unconditional (blank prompt)
- Negative Prompt
- blurry, cropped, ugly

- Prompt
- A photo-realistic image of a cat
- Negative Prompt
- blurry, cropped, ugly

- Prompt
- prompt not found (2)
- Negative Prompt
- blurry, cropped, ugly

- Prompt
- prompt not found (3)
- Negative Prompt
- blurry, cropped, ugly
The text encoder was not trained. You may reuse the base model text encoder for inference.
Training settings
Training epochs: 49
Training steps: 400
Learning rate: 0.0001
- Learning rate schedule: constant
- Warmup steps: 0
Max grad value: 2.0
Effective batch size: 3
- Micro-batch size: 1
- Gradient accumulation steps: 1
- Number of GPUs: 3
Gradient checkpointing: True
Prediction type: epsilon (extra parameters=['training_scheduler_timestep_spacing=trailing', 'inference_scheduler_timestep_spacing=trailing'])
Optimizer: bnb-lion8bit
Trainable parameter precision: Pure BF16
Base model precision:
no_change
Caption dropout probability: 0.1%
LoRA Rank: 128
LoRA Alpha: 128.0
LoRA Dropout: 0.1
LoRA initialisation style: default
Datasets
antelope-data
- Repeats: 0
- Total number of images: ~24
- Total number of aspect buckets: 1
- Resolution: 1.048576 megapixels
- Cropped: True
- Crop style: center
- Crop aspect: square
- Used for regularisation data: No
Inference
import torch
from diffusers import DiffusionPipeline
model_id = 'stabilityai/stable-diffusion-xl-base-1.0'
adapter_id = 'bghira/simpletuner-controlnet-sdxl-lora-test'
pipeline = DiffusionPipeline.from_pretrained(model_id, torch_dtype=torch.bfloat16) # loading directly in bf16
pipeline.load_lora_weights(adapter_id)
prompt = "A photo-realistic image of a cat"
negative_prompt = 'blurry, cropped, ugly'
## Optional: quantise the model to save on vram.
## Note: The model was not quantised during training, so it is not necessary to quantise it during inference time.
#from optimum.quanto import quantize, freeze, qint8
#quantize(pipeline.unet, weights=qint8)
#freeze(pipeline.unet)
pipeline.to('cuda' if torch.cuda.is_available() else 'mps' if torch.backends.mps.is_available() else 'cpu') # the pipeline is already in its target precision level
model_output = pipeline(
prompt=prompt,
negative_prompt=negative_prompt,
num_inference_steps=20,
generator=torch.Generator(device='cuda' if torch.cuda.is_available() else 'mps' if torch.backends.mps.is_available() else 'cpu').manual_seed(42),
width=1024,
height=1024,
guidance_scale=4.2,
guidance_rescale=0.0,
).images[0]
model_output.save("output.png", format="PNG")