readme: # Stable Diffusion web UI
  A browser interface based on Gradio library for Stable Diffusion.
  
  ![](txt2img_Screenshot.png)
  
  Check the [custom scripts](https://github.com/AUTOMATIC1111/stable-diffusion-webui/wiki/Custom-Scripts) wiki page for extra scripts developed by users.
  
  ## Features
  [Detailed feature showcase with images](https://github.com/AUTOMATIC1111/stable-diffusion-webui/wiki/Features):
  - Original txt2img and img2img modes
  - One click install and run script (but you still must install python and git)
  - Outpainting
  - Inpainting
  - Color Sketch
  - Prompt Matrix
  - Stable Diffusion Upscale
  - Attention, specify parts of text that the model should pay more attention to
      - a man in a ((tuxedo)) - will pay more attention to tuxedo
      - a man in a (tuxedo:1.21) - alternative syntax
      - select text and press ctrl+up or ctrl+down to automatically adjust attention to selected text (code contributed by anonymous user)
  - Loopback, run img2img processing multiple times
  - X/Y plot, a way to draw a 2 dimensional plot of images with different parameters
  - Textual Inversion
      - have as many embeddings as you want and use any names you like for them
      - use multiple embeddings with different numbers of vectors per token
      - works with half precision floating point numbers
      - train embeddings on 8GB (also reports of 6GB working)
  - Extras tab with:
      - GFPGAN, neural network that fixes faces
      - CodeFormer, face restoration tool as an alternative to GFPGAN
      - RealESRGAN, neural network upscaler
      - ESRGAN, neural network upscaler with a lot of third party models
      - SwinIR and Swin2SR([see here](https://github.com/AUTOMATIC1111/stable-diffusion-webui/pull/2092)), neural network upscalers
      - LDSR, Latent diffusion super resolution upscaling
  - Resizing aspect ratio options
  - Sampling method selection
      - Adjust sampler eta values (noise multiplier)
      - More advanced noise setting options
  - Interrupt processing at any time
  - 4GB video card support (also reports of 2GB working)
  - Correct seeds for batches
  - Live prompt token length validation
  - Generation parameters
       - parameters you used to generate images are saved with that image
       - in PNG chunks for PNG, in EXIF for JPEG
       - can drag the image to PNG info tab to restore generation parameters and automatically copy them into UI
       - can be disabled in settings
       - drag and drop an image/text-parameters to promptbox
  - Read Generation Parameters Button, loads parameters in promptbox to UI
  - Settings page
  - Running arbitrary python code from UI (must run with --allow-code to enable)
  - Mouseover hints for most UI elements
  - Possible to change defaults/mix/max/step values for UI elements via text config
  - Random artist button
  - Tiling support, a checkbox to create images that can be tiled like textures
  - Progress bar and live image generation preview
  - Negative prompt, an extra text field that allows you to list what you don't want to see in generated image
  - Styles, a way to save part of prompt and easily apply them via dropdown later
  - Variations, a way to generate same image but with tiny differences
  - Seed resizing, a way to generate same image but at slightly different resolution
  - CLIP interrogator, a button that tries to guess prompt from an image
  - Prompt Editing, a way to change prompt mid-generation, say to start making a watermelon and switch to anime girl midway
  - Batch Processing, process a group of files using img2img
  - Img2img Alternative, reverse Euler method of cross attention control
  - Highres Fix, a convenience option to produce high resolution pictures in one click without usual distortions
  - Reloading checkpoints on the fly
  - Checkpoint Merger, a tab that allows you to merge up to 3 checkpoints into one
  - [Custom scripts](https://github.com/AUTOMATIC1111/stable-diffusion-webui/wiki/Custom-Scripts) with many extensions from community
  - [Composable-Diffusion](https://energy-based-model.github.io/Compositional-Visual-Generation-with-Composable-Diffusion-Models/), a way to use multiple prompts at once
       - separate prompts using uppercase `AND`
       - also supports weights for prompts: `a cat :1.2 AND a dog AND a penguin :2.2`
  - No token limit for prompts (original stable diffusion lets you use up to 75 tokens)
  - DeepDanbooru integration, creates danbooru style tags for anime prompts
  - [xformers](https://github.com/AUTOMATIC1111/stable-diffusion-webui/wiki/Xformers), major speed increase for select cards: (add --xformers to commandline args)
  - via extension: [History tab](https://github.com/yfszzx/stable-diffusion-webui-images-browser): view, direct and delete images conveniently within the UI
  - Generate forever option
  - Training tab
       - hypernetworks and embeddings options
       - Preprocessing images: cropping, mirroring, autotagging using BLIP or deepdanbooru (for anime)
  - Clip skip
  - Use Hypernetworks
  - Use VAEs
  - Estimated completion time in progress bar
  - API
  - Support for dedicated [inpainting model](h