Flow diversion for cerebral aneurysms or stabilizing the hemodynamics? I'm doing research of flow diverter for cerebral aneurysms applications and I'm wondering the reason behind stent placement underneath of cerebral aneurysms is to divert the flow or stabilize the hemodynamics (e.g. reduce wall shear stress or blood flow velocity in the aneurysmal region)? Let's say people do it to divert the flow and prevent blood to enter the aneurysmal region. That looks fine for side wall aneurysms when there is no smaller arterioles or branches connected to the aneurysm, but for terminal aneurysm where the flow diverter may block the blood flow into smaller arterioles, is it a appropriate case to deploy stent? I mean if we prevent the blood to enter some part of smaller arterioles how brain could maintain its functionality due to its need for oxygen. I appreciate if someone could explain the logic behind the flow diverter placement in the cerebral aneurysms.

I don't quite understand the dichotomy you present between diverting flow and stabilizing hemodynamics. The former is a direct outcome and the latter is a secondary outcome: both are goals.

The overall goal is to isolate the blood in the aneurysm and allow it to clot and heal over time. Embolic coils are also used to assist the process, in which case the flow diverter helps to keep the coil in place and the resulting embolism out of the primary circulation.

One would not typically deploy a flow diverter over branches that one wishes to remain patent, although patency can be maintained in smaller vessels. In many neurovascular situations, collateral circulation can make it acceptable to occlude a branch, however.

You may be confusing plain stents with flow diverters; a plain stent has as few struts as possible while maintaining structural support (though high flow through the struts can cause turbulence and is something you want to avoid if possible), and creative stent placement techniques including deployment of multiple stents and use of balloons to break open strut configurations can be used. These do not block flow per se, but can help keep coils in place and maintain the desired structure of the vessels involved.

I've included several references that I quickly gathered; these may not be the most ideal references but should give you a look at some of the approaches used. I would strongly suggest you look through the rest of the literature using search terms from your question to get an idea of how interventionalists approach different situations.

I'll add that I am not a physician and am not up-to-date with the latest guidelines or standard of care, but in my experience with the literature a lot of these decisions are up to the experience of the physicians involved and based on a lot of case study; randomized trials aren't really possible to cover all the anatomical variety that is experienced.

Burkhardt, J. K., Riina, H. A., Tanweer, O., Shirani, P., Raz, E., Shapiro, M., & Nelson, P. K. (2018). Flow diversion and microvascular plug occlusion for the treatment of a complex unruptured basilar/superior cerebellar artery aneurysm: case report. Journal of neurosurgery, 1-6.

Chow, M. M., Woo, H. H., Masaryk, T. J., & Rasmussen, P. A. (2004). A novel endovascular treatment of a wide-necked basilar apex aneurysm by using a Y-configuration, double-stent technique. American Journal of Neuroradiology, 25(3), 509-512.

Miyachi, S., Matsubara, N., Izumi, T., Asai, T., Yamanouchi, T., Ota, K., ... & Wakabayashi, T. (2013). Stent/balloon combination assist technique for wide-necked basilar terminal aneurysms. Interventional Neuroradiology, 19(3), 299-305.

Puffer, R. C., Kallmes, D. F., Cloft, H. J., & Lanzino, G. (2012). Patency of the ophthalmic artery after flow diversion treatment of paraclinoid aneurysms. Journal of neurosurgery, 116(4), 892-896.

Siddiqui, M. A., Bhattacharya, J. J., Lindsay, K. W., & Jenkins, S. (2009). Horizontal stent-assisted coil embolisation of wide-necked intracranial aneurysms with the Enterprise stent—a case series with early angiographic follow-up. Neuroradiology, 51(6), 411-418.

Wong, G. K., Kwan, M. C., Ng, R. Y., Simon, C. H., & Poon, W. S. (2011). Flow diverters for treatment of intracranial aneurysms: current status and ongoing clinical trials. Journal of Clinical Neuroscience, 18(6), 737-740.


Osmosis in red blood cells and bacteria This is a question from an exam in my biology course.

Bacterial cells and human red blood cells were inserted into one solution. Upon testing one hour later the blood cells exploded, while the bacterial cells stayed intact. What answer explains the findings:

A. The bacteria cell walls prevented their explosion.

B. The blood cell walls didn't prevent their explosion.

C. solute concentration in the solution is equal to the blood cells and bacteria solute concentration.

D. solute concentration in the solution is higher than the blood cells solute concentration.

My lecturer said the right answer is B. I thought the right answer is A.
What would be the correct answer and why?

To my understanding, red blood cells don't have a cell wall, but rather a membrane, which is why in my opinion you can't say that it prevented the explosion if it doesn't exist. On the other hand, bacteria does have cell wall, that's why I thought this answer is the correct one.

Absent a very specific learning objective, if (B) is correct, this is a poor question. In a high-school or undergraduate introductory biology course, I would say:

  (A) The bacteria cell walls prevented their explosion

Would be the correct answer, because of exactly your reasoning. This is what you're taught in introductory biology. You're taught that animal cells don't have cell walls. Bacteria and plant cells do have cell walls.

See The Cell

  animal cells are not surrounded by cell walls

Under this concept of a cell wall, red blood cells certainly don't have them.

However, red cell osmotic fragility testing comes out of a tradition that uses the phrase cell wall strength to describe the response of a red cell to osmotic stress. See Orcutt

  It appeared logical to us that a model based on a distribution of cellular wall strengths ought to reproduce the experimentally observed variation of degree of hemolysis as a function of osmotic stress.

All the same, if you were operating from that perspective, how would you distinguish between (A) and (B)? Here, we're not even taking into account the fact the primary mechanism for responding to osmotic stress in most animal cells, including the human erythrocyte, is regulation of ion flow, not the pressure of a rigid extracellular matrix. If the solution contained ouabain, which certainly does lyse human RBCs, would it be best to characterize the problem as "The blood cell walls didn't prevent their explosion"?

I would ask your lecturer to clarify the learning objective behind the question. As an interpersonal note, I'd point out that this is very different from asking him or her to justify the correct answer. Some (many...) teachers are oddly sensitive to any suggestion that they need to justify their expertise, especially when it comes to exam questions. If you make it a question about the learning objective, this should help you avoid triggering that particular sensitivity.


What distinguishes "coding" from "noncoding" DNA? I've been reading a bit about "junk DNA" and how much of our genome consists of this "non coding DNA" in comparison to "coding DNA".

I'm just an interested layperson but I thought all combinations of three base pairs encoded one amino acid, with some amino acids being encoded by more than one combination of base pairs.

But if that were true then all of our DNA would encode something.

Or if only a tiny percentage of our DNA is "coding" that would mean that the vast majority of possible combinations of three base pairs don't represent any amino acid.

Or it could mean that there are a small number of "meaningless" combinations of three base pairs, but that those combinations are vastly overrepresented in our genome.

Which is correct? What am I missing?

You are correct in thinking that any sequence of bases corresponds, via the genetic code, to a sequence of amino acids. However not all stretches of DNA are actually transcribed into mRNA for translation into proteins. For this to happen the stretch of DNA requires (DNA-encoded) elements to promote and regulate the transcription and translation processes, and this, very broadly, is what defines a gene: a segment of DNA which has the required components to direct the synthesis of a protein (or in some cases an RNA that will not be translated into a protein). The segments of DNA that correspond to mRNA for proteins and to other RNA molecules are referred to as coding sequences.

Now, I've omitted lots of details here: some genes encode RNA molecules, such as ribosomal RNAs, which are not translated into proteins; eukaryotic genes include stretches of sequence (called introns) that are spliced out of the transcript before the mRNA>protein step. The original definition (if there ever was a definition) of "junk" DNA included these introns, as well as regions of DNA lying outside the coding sequences. We now know that there is useful information stored in much of this DNA, even though it doesn't code for anything directly via the genetic code.


Cut strand of hair longevity? Quick and useless question : what can be the longevity of a cut strand of hair, assuming we kept it in (good) normal conditions ? Any idea ?

In order to keep the hair from decomposing, my first assumption is you'd need it in an environment that inhibits decomposers, e.g. dry. We can see in the mummification process that hair may be preserved for hundreds of years. Hair is also composed largely of keratin, however, and in my own search of the literature: (1) we need microbial (bacteria, fungi) keratinases to decompose the keratin in hair [source], and (2) this process is slow. There's a really old article here, from the abstract:

  Keratin, prepared from horn meal, was added to moist field and garden soil and allowed to decompose in the laboratory. The keratin was found to undergo a decomposition resulting in a slow, but steady accumulation of ammonia and nitrate. 35–40 per cent, of its nitrogen was transformed into nitrate after 120 days.

More recent numbers to shed light on a rate of degradation for keratin are either obscure, or untested. This is more so for hair specifically as I couldn't find much more.


Tryptophan Operon Working Mechanism The Trp operon has 2 methods to check for presence of tryptophan, one indicating the absence of free Trp and the other the absence of charged Trp tRNA. If both methods showed absence of Trp there is tryptophan synthesis. However, the synthesis of tryptophan involves 5 enzymes to be synthesized, and usually we find all the amino acids in a protein such as an enzyme. Therefore, there is high probability that Trp is found in these 5 enzymes. Where does that Trp come from if there is neither free tryptophan nor charged tRNA for Trp?

You're right that if there was really a complete lack of Trp, then the cell wouldn't be able to make the enzymes. But the cell doesn't have to wait until [Trp] drops all the way to zero -- it can induce expression of the operon when the concentration gets sufficiently low (in practice there's some continous relationship between [Trp] and expression). Note that there's always going to be a little bit of Trp around as there are constantly proteins being degraded (and resynthesized).

But you're on to something -- in fact there is some bias towards lower usage of a particular amino acid in enzymes that are responsible for synthesizing that amino acid. Here's the paper: https://www.ncbi.nlm.nih.gov/pmc/articles/PMC1839009/ although the effect isn't that high for Tryptophan.


What species is this little white bird? This bird was found in my backyard in the north east of Italy. It's about 5cm long and can fly as high as 1.5m (so far).


What is the relation between plasmid concentration and mRNA levels? Suppose a simple synthetic construct, consisting of a constitutive promoter and a single gene:

One of the simplest ways to model GFP transcription is to use an ODE:

$\frac{d [GFP_{mRNA}]}{dt} = a - b{\cdot}[GFP_{mRNA}]$

where $a$ is GFP transcription rate and $b$ is GFP mRNA degradation rate (both constants). Normally, we assume $a>>b$.

Suppose we wish to account for plasmid concentration with the value of transcription rate $a$ - e.g. with higher plasmid concentration, the ratio $a/b$ should increase as the mRNA saturation levels are expected to be reached faster. In other words, suppose we transfect 2 individual constructs depicted above, 10 ng of one and 30 ng of the other - how should this difference in concentration impact the rate constants values, assuming the constructs are otherwise identical?

Assuming this, approximately how does gene transcription rate $a$ change as a function of the amount of transfected plasmid containing the above construct? Ideally, I'd be interested in knowing this for HEK 293 cells, but any other decent estimation is acceptable. One simple option is to simply assume linear relation, but I think that's too simplistic: e.g. if we say $a=1$ at 10 ng plasmid concentration, then $a=3$ at 30 ng plasmid concentration. Note that I'm not necessarily limiting myself to ODE models, just looking for some kind of relation.

I have so far not been able to find anything useful in papers I've examined, so any help would be appreciated.

The amount of transfected plasmid does not correlate at all with the protein expression level. After transfction, usually each cell is going to get and keep only one copy of the plasmid. Once the plasmid is in the cell, it will be replicated and the cell will contain X copies of it, depending on the plasmid copy number. In general, plasmids with low copy numbers correlate with lower protein expression compared to plasmids with high copy number. But there are also exceptions:

  Expression of genes under control of EF1α promoter appears to correlate with plasmid copy number.  In contrast, expression driven by the more traditional CMV immediate early promoter appears less plasmid copy number dose-responsive (ref).


Why is Electrical Cardioversion contraindicated in AF and Digitalis Poisoning? I am thinking why electrical cardioversion is contraindicated in atrial fibrillation with digitalis toxicity/poisoning.

Cardiac digitalis is also called digoxin and cardiac glycoside.
It is antiarrhythmic agent against supraventricular arrhythmias.
It is contraindicated in Atrial Fibrillation with digitalis toxicity.
The Pubmed's article about the topic is a stub without describing its Pharmacology.

Digitalis toxicity is a complication of digitalis therapy, here used for atrial fibrillation.
This may be due to kidney failure when the medication accumulates in the body.
Its treatment is activated charcoal and gastric lavage.
However, assume that these methods do not help in treatment or that we in acute digitalis toxicity.

I think the reason why we cannot recently do electrical cardioversion in digitalis poisoning is the high probability of repolarizing incidences so the likelihood to stimulate heart at the right moment is very small.
This is because digitalis poisoning is associated with cardiac disturbances: ventricular tachycardia, ventricular fibrillation, SA block and AV block.
Without knowing more precisely about the complication, I think the electrical cardioversion has high risk of failure.

Why is electrical cardioversion contraindicated in Digitalis Poisoning of AF treatment?

The mechanism of action of digitalis-like preparations based on inhibition of myocardial NA+/K+-ATPase as well as of cardiac conducting tissue. The inhibition of latter (conducting) is 2.5 greater than those of former (myocardial). This is why the main feature of digitalis intoxication are arrhythmias due to ectopic electric impulses and spontaneous depolarization (there is overloading of conducting system by Calcium as well as excess amounts of Noradrenaline released from sympathetic endings).

These metabolic changes in a case of digitalis intoxication could lead to irreversible ventricular arrhythmias after cardioversion - there are sporadic reports about this phenomenon. Based on these reports, it was thought that the use of digitalis during electrical cardioversion would provoke ventricular arrhythmias in patient even with clinically desirable concentration of digitalis.

In a study in patients with therapeutic nontoxic digitalis levels determined that these patients may safely undergo electrical cardioversion without reducing or withholding the dose of digitalis. Thus, the cardioversion is not contraindicated in a patient without signs of digitalis toxicity.

Although, the exact mechanism of the ventricular arrhythmias after cardioversion is not well understood, it is clear that mentioned above changes in electrolytes concentrations with electrical stimulation lead to premature beats overtaking restored sinoatrial pace. The rate of ventricular arrhythmias is in direct correlation with the level of the energy used for the cardioversion (this is why there is "pre-cardioversion" with low-energy pulses). Until these metabolic changes are not corrected, the resulting ventricular arrhythmias can be irreversible. This is why, digitalis intoxication is relative contraindication to electric cardioversion of atrial fibrillation/flatter.

Excerpt from the article mentioned:


Exact definition of homoplasy? I am a poor mathematician looking into the Sars-Cov-2 genome. Is there a proper definition of the word homoplasy? Is it just mutations that happened independently or does there have to be indications that the mutations give an improvement in fitness?

Homoplasy just refers to the fact that the mutation happened independently. While several sources only discuss homoplasy as a phenotypic characteristic, and in contrast to homology (and thus suggest an adaptive implication), in general,  there is no assumption that the mutation gave increase in fitness (see this paper).


Does posture matter while drinking water? I have heard a lot of people (where I am from, India) say that one must drink water sitting, not standing. Recently, I also came across a video that said the same thing (link). The reason given is that when you drink water standing, it goes to the colon too quickly, whereas while sitting, the water is better absorbed by the organs. Is this true? Is there a biological basis for why I should drink water sitting?

I tried to find literature on this matter, but apparently there isn't a whole lot. I found this paper where they tried to see the water passage rate in the GI tract of rabbits, mentioning how the water goes from stomach to cecum to proximal large colon, but obviously did not study posture.

Note: I think this question belongs here better than Skeptics.SE as I am looking for a biological explanation

The stomach is equipped with the pyloric valve to regulate transpyloric flow to the duodenum.
Both cold and warm drinks stimulate a pattern of motility associated with retardation of transpyloric flow.

The pyloric valve is responsible for acidifying water that you drink, and if the body's guard against amoeba's and dysentery depended on body position, it would be low survival performance, given that sedentary posture is associated with poorer health outcomes, including an increased risk of type 2 diabetes. You will benefit from minimizing time spent sitting each day, and from breaking up periods of time spent being sedentary, as often as possible. It's aperture at rest is 3mm, so it's of minimal difficulty for the stomach to hold water as required for optimal health. It's highly likely that water is kept for enough time to expose it to stomach acids, regardless of body position, and then it's transmitted.
Water doesn't flow straight through the stomach like a straight tube when it's empty.
https://www.ncbi.nlm.nih.gov/pmc/articles/PMC1382811/pdf/gut00528-0043.pdf


Need some good resources to learn about Protein function and structure I'm taking a course on biochemistry at edx.
Since I'm a computer science student, I'm having some trouble in understanding many biochemical concepts. While the first module was just fine, I found the second module really hard to understand. I have problems to learn and understand about amino acids and proteins, so I would like to study apart from some good resources.

Which is a good resource to learn about amino acids and macromolecules?
Also, please let me know if this book is up to date?

Given your background (not a biologist or chemist) you would probably find the introductory material in a biochemistry "lite" textbook more accessible/useful than a hardcore text for specialists.  As a co-author of a biochemistry textbook, I can tell you that there are essentially three different classes, or types of texts.

Comprehensive textbooks, of which I think Voet & Voet is an outstanding example, are designed for biochemistry majors who will be taking a full year, in-depth course, typically in the 3rd year of their undergraduate studies.
The general introductory texts are usually thinner, and are designed to whip through the same material in half a year--aimed at life-science students in other majors. Some of these books that I described as "lite" are further targeted to, for example, nursing, medical, or dental school students.
The third category of text tries to be multi-purpose, and appeal to both markets. The current versions of Stryer and Lehninger fall into this class (they are both excellent in my opinion); the instructors can cherrypick the readings if they want to use the book for a half-year course.

Voet, Voet, and Pratt (mentioned in the comments above) is a very good diluted version of Voet & Voet. As you can also tell from the comments, different texts have different flavors, and so which one is best depends on the reader.

It sounds like a more general "concept" based book might be useful for you (to provide background and context). The very thick "Molecular Biology of the Cell" by Alberts et al., has a lighter cousin named Essential Cell Biology which would be well worth looking into.

Keep in mind that when tackling a university-level biochemistry course the instructor has to assume that you have passed university level introductory courses in general chemistry and in organic chemistry. Biochemistry majors typically complete full year courses in both these subjects, as well as a watered-down version of physical chemistry.


Do eukaryote cells contain DNA that isn't part of a chromosome or located in the mitochondria? I specify eukaryote in the title, but I'm also interested if this question isn't applicable to eukaryote cells in general but is to humans. I was reading "RNA-seq: An assessment of technical reproducibility and comparison with gene expression arrays" (John Marioni 2008).

In the results it states

"By these criteria, 40% of reads mapped uniquely to a genomic location, and of these, 65% mapped to autosomal or sex chromosomes (the remainder mapped almost exclusively to mitochondrial DNA)."

I couldn't help but notice the "almost exclusively to mitochondrial DNA". Almost exclusively? Can DNA be found in places other than chromosones or mitochondria? Perhaps I'm interpreting the sentence wrong. Any pointers would be appreciated

Thanks

In plants, chloroplasts and other plastids contain DNA, but I suppose you are more interested in humans. Quoting from wikipedia,

  In many cells cytoplasmic DNA is also found, which is different from
  nuclear DNA, both in methylation levels (cytoplasmic has less), and in
  sequence. EccDNA or extrachromosomal circular DNA is present in all
  eukaryotic cells, derived from genomic DNA and consists of repetitive
  sequences of DNA found in both coding and non-coding regions of
  chromosomes. EccDNA can vary in size from less than 2000 more than
  20,000 base pairs. In animals, eccDNA molecules have been shown to
  contain repetitive sequences that are seen in satellite DNA, 5S
  ribosomal DNA and telomere DNA. The function of eccDNA has not been widely studied, but it has been proposed that the production of elements of eccDNA from genomic DNA sequences adds to the plasticity of the eukaryotic genome and can influence genome stability, cell aging and the evolution of chromosomes


How do plants store energy by potassium battery? Plants have a so-called potassium battery system for energy storage/conversation. Wikipedia only has a reference to a technical article. Can someone explain it in a layman's language (or refer to a non-technical article)?

I'll give it a layman's attempt in sort of an "ELI5" approach. I hope the analogy isn't too simplistic, and please don't take it too literally.

How do plants use K+?

Plants control K+ concentrations and do work within the plant by pumping K+ between different compartments. Because water tends to follow solute concentrations (osmosis), plants can do things like suck in water simply by moving potassium around.

(note: animal cells do something similar: this is how your nervous system, kidneys, digestive tract, and other cells/organs perform their functions, but using different ions or combinations of ions)

Analogy to some human activities

Moving all this K+ around takes a lot of energy. Creating a concentration gradient is a lot like pumping water uphill: you put some energy in, but you are storing that energy as potential energy for the water flowing back downhill, or the ions flowing back against their concentration gradient (or water flowing the other way). But normally, plants use this as a way to spend energy to do useful things, rather than a way to store it. It's like they are pumping the water up the hill and then floating useful ores or lumber on a boat as it flows back down the hill.

What is actually happening that Wikipedia describes as a plant "K+ battery"

The article cited on Wikipedia, though, Gajdanowicz et al 2011, shows that sometimes plants take advantage of the mobility of potassium ions to assist the H+ ATPase that is another form of energy storage. Plants use H+ (proton) concentration gradients similarly to how they use K+ concentration gradients, but they are used to power different processes in the plant. Sort of like two different types of batteries.

Return to the analogy, and why this is useful to the plant

Normally both H+ gradients and K+ gradients are powered by ATP, which is indirectly produced via photosynthesis: think of it like the plant's coal-fired power plant. Converting K+ concentration gradients to H+ concentration gradients would be like using the water you pumped up the hill to power a hydroelectric plant that is charging a car battery. The benefit for the plant is that, because K+ can be easily moved around, the plant is able to power those hydroelectric plants anywhere around the hill, and that can be advantageous if a particular part of the plant is low on other sources of energy.


Do patella bones form in people born with the inability to walk? The patella is a sesamoid bone that typically doesn't completely from and ossify until ~3-6 years of age (e.g., Source). My long-standing understanding (supported by a claim in Saladin's college A&P textbook1 sitting on my shelf) was that sesamoid bones form "in response to stress."
Well, I had a student in my anatomy class this morning ask me the titular question about people born with leg paralyses or other congenital issues preventing them form walking. To their point, if such congenitally paralyzed people never walk and therefore never place stress on the knee via the quad tendon, would they ever form a sesamoid bone (i.e., the patella) in the knee?
I've done some digging and found evidence of "nail patella syndrome" (e.g., here) which results in people not forming a complete patella. However, I was unable to quickly narrow down any sources discussing lack of patellar growth due to congenital leg paralysis.

1.  Saladin, K. S. 2015. Anatomy & Physiology: The Unity of form and function. 7th edition. McGrawHill, Education, New York, NY

In mice, patella would form, but not separate from femur.
On the development of the patella:

Previous studies have suggested a central role for mechanical load in
the initiation of sesamoid development.
Our finding that the patella initially forms as part of the femur and
is later separated by joint formation led us to revisit the issue of
mechanical contribution to patella development. To achieve this, we
studied patella development in mutant embryos devoid of contracting
muscles, as a result of the naturally occurring autosomal recessive
mutation muscular dysgenesis (mdg; Cacna1s). These mice lack
excitation-contraction coupling and, consequently, their skeletal
muscles lack contractility.
Histological analysis of sagittal sections of hindlimbs identified at
E13.5, in both mutant and control littermates, the existence of
patella progenitors, which were part of the femur. However, at E14.5,
whereas in the control we observed joint formation between the patella
and the femur, in mdg mutants this process was not detected.
Consequently, at E16.5 the mutant patella was still part of the femur.
These results suggest that, although muscle contraction does not
affect patella initiation, it plays a crucial role later during the
separation phase.

I know, you asked for human, but it is the only relevant information I found. However, while sources mention a role of physical stress in the calcification of sesamoid bones, the preceding formation of cartilage seems unaffected. Now, given the information that ossification occurs in mice without physical stress, I personally would conclude that the human patella would also calcify in paralyzed persons. Strictly, the question whether the patella forms, should be answered, if a non-detached patella counts.


Human can't digest cellulose.....but the length of small intestine is approximately same in humans and herbivores My question is why herbivores are able to digest cellulose and why can't we??

Depends on what you mean by "digest". Herbivores aren't able to digest grass on their own, they keep bacteria in their guts to do that. Grass-eating herbivores also tend to have long digestive cycles and some have several stomach chambers.

In short, cellulose is a pain to digest.


Questions on adding a protein to a DNA library Two questions regarding finding the DNA sequence of a amino acid sequence (AA):

1) If you are able to find out the mRNA sequence of an AA, then don't you automatically know the DNA sequence?

2) When determining the mRNA sequence of an AA, how do you account for the various permutations in the genetic code for a given AA. Ex. Leucine is CUU and CUC. Or does it not matter?

For 1, that is how retroviruses such as HIV work, and of which scientists have utilised to use in gene therapy. It requires a process called reverse transcription.
However, this would only show you the gene the mRNA was from, but no information about its location, or if it is repeated elsewhere.

As for 2, this is to do with the fact that one base cannot code for one amino acid; as there are 20 common amino acids used to make proteins. With 2 bases, you are only getting 16 possible codes from the various base pairs you can create. However, with 3 bases, 64 possible combinations are possible, more than enough to suffice for the 20 needed. This is called the triplet code;
and therefore, there are more possible combinations of the base pairs in DNA than there are amino acids to be made from them. The 'left over' combinations can code for the same amino acids as other combinations, making it the code degenerate. Therefore, it doesn't matter, as the combination is always read in 3s (non-overlapping), and if it codes for a particular amino acid, then that amino acid will be produced at the ribosomes when its corresponding code is read.

I hope this increases your understanding of the topic.

http://www.genetherapynet.com/viral-vector/retroviruses.html (Question 1)
http://www.chemguide.co.uk/organicprops/aminoacids/dna4.html (Question 2)


Cattail seeds spring revitalization There's a genus of plants known as Typha, or cattail.
They have seeds like of dandelion, that form thick heads. Such a head may be cut off from the plant and be used for, say, decorative purpose, but when the spring comes, seeds will spread out their 'hairs' and fly away. It looks like an absolutely dry piece of plant - no life in it - suddenly gets alive.

I suggest that it may be a daytime prolongation of a  that causes it to spread seeds, but I wonder if anyone knows what sort of cellular/biochemical mechanics is at work here.

According to LeRoy G. Holm's World Weeds: Natural Histories and Distribution:

  Seeds are protected from germinating in the stalk by the swelling of the pistillodes that are distributed throughout the inflorescence. In humid conditions, they swell and keep the spike intact. They can swell 15 to 20% in 15 min in 100% relative humidity. Upon drying, the pistillodes shrivel, air again reaches the seed, and the perianth hairs spread. Their collective pressure causes the inflorescence to burst, releasing the fruits to the air. If they fall on water, the hairs close and the pericarp releases the seed which then sinks.

So it's all about the humidity!


Why shouldn't there be any air between lungs and ribcage? My OL biology textbook says that one of the functions of the pleural fluid between lungs and pleural membrane is that it ensures no air between lungs and rib cage.

But I wonder why air there would cause any problems. I was told that it'd build up there and exert pressure on the lungs but why would it build up there in first place, couldn't it just go out through the mouth or nose? And in case that's really the answer, how does it manage to prevent any oxygen or other gases from getting into the fluid, after all, the pleural fluid touches the pleural membranes cells which require oxygen, won't oxygen diffuse from around them into the fluid?

Pleural fluid is secreted by the pleural membrane and therefore it is not made up of living cells. Usually a trauma to the chest, such as a stab wound, would allow air to enter the thoracic cavity, called pneumothorax.

The problem caused by air building up in the thoracic cavity is due to the movement of the diaphragm and intercostal muscles cannot create a low pressure relative to that of within the lungs pressure, and air rushes into the thoracic cavity instead of the lung to equalizes the pressure. Note that oxygen that enters the thoracic cavity is not usable by the individual, and only the oxygen that enters the lung can. Also, air in the thoracic cavity is not connected to the nose or mouth, or else only little air would enter the lung when aspirating.


Snail species identification from the United Kingdom Can anybody identify the species of this snail?

From this website, it looks like a Common Snail.

Could it be the dextral type of Cornu aspersum?

These snails have been found in a garden, have survived a month on lettuce in a home environment, retreat in their shell sometimes. Their shell measures about 2 cm diameter.

This snail is Cornu aspersum its common name is garden snail, it belongs to Helicidae family click here for more information


Protein Structure Parameters I'm wondering about the minimal set of parameters necessary to define a protein's structure.  My understanding is that the backbone geometry is defined by the phi and psi angles (torsion angles), and then from this template the side-chain rotamers are defined.  Is this sufficient to yield a geometric structure, or are there other factors involved?

On a related note, how exactly are two protein structures compared? I often see references to percent similarity, but what does this really mean?  In the Protein Data Bank, atomic coordinates are given, but my intuition is that coordinates are only important relative to all the other amino acids in the chain.  In other words, absolute coordinate comparison would seem to be uninformative.

It depends on how accurate you want the description to be. Ideal bond lengths are normally assumed, so the structure can be described using φ, ψ, ω (for the backbone) and χ1-χ5 (for the sidechain) angles. This is what is used in most cases. However, the positions of the hydrogens in the sidechains are not uniquely defined by this description. This also does not describe hydrogen bonds and disulfide bridges, but those could be deduced from the relative positions.

You can get more information here: http://kinemage.biochem.duke.edu/teaching/anatax/html/anatax.1b.html

This description should be enough for most applications, though.


What controls gut motility? I have two different papers. One claims that gut motility is reduced by stimulation of the Opioid κ and δ receptors. The receptors are activated by Morphine and certain derivatives, specifically Codeine, Fentanyl, Hydromophone, Methadone, and Oxycodone.

The second paper claims that gut motility is reduced by blocking the 5-HTb3 receptors.  Antagonists or competitors to the 5-HTb3 receptors being Morphine and certain derivatives.

Are these two different theories, two different action mechanisms, or am I missing something?

When discussing the control of gut motility there is more to mention than the use of serotonergics and opioids - a number of peptide and nonpeptide neurotransmitters are important. Somatostatin and nitric oxide are two examples which each happen be inhibitory of intestinal motor function, but by two completely different mechanisms. A bit of understanding in cell biology, specifically cell signalling would be helpful to fully grasp this, but the former acts by inhibiting the release of Acetylcholine and Substance P while the latter increases levels of cAMP/cGMP. And don't forget about the basics when it comes to motor control - the cholinergics and adrenergics (including dopamine) are logical opposites.

References
Leung, P.S. 2014, The Gastrointestinal System: Gastrointestinal, Nutritional and Hepatobiliary Physiology, Springer, 19/03/2018, http://www.springer.com/cda/content/document/cda_downloaddocument/9789401787703-c2.pdf?SGWID=0-0-45-1491980-p176563984


How much heat can a human sustain? For how long can a human sustain a temperature of 180 degrees Fahrenheit (82oC) without damage in a confined place?

For instance, suppose a person is sitting in a steam-filled room, without external air flow (but with enough oxygen) - how long could one resist the heat before permanent physical damage would occur?

Short answer
Water of 180 oF (82 oC) causes immediate scalding (thermal burn wounds).

Background
The severity of a scald injury depends on the temperature to which the skin is exposed, and for how long. Residential water heaters warm up tap water typically to 120 oF (48 oC). At this temperature, the skin of adults requires an average of five minutes of exposure for a full thickness burn to occur. When the temperature of a hot liquid is increased to 140 oF (60 oC)) it takes only five seconds or less for a serious burn to occur. Coffee, tea, hot chocolate and other hot beverages are usually served at 160 to 180 oF (71 - 82 oC), resulting in almost instantaneous burns that will require surgery (American Burn Association).

How these numbers apply to steam I am not sure, but I think they will be comparable, given that steam is essentially made up of a dense mass of water droplets. Further, the amount of steam obviously matters. Consider these numbers: The temperature in a sauna typically ranges from 160-200 oF (71 - 93 oC) with a low level of humidity (ranging from 5 - 30%). Steam rooms provide moist heat from a water-filled generator pumping steam into the enclosed room. The temperature in a steam room typically ranges from 110-114 degrees Fahrenheit (43 - 46 oC) with a humidity level of 100%.

There is anecdotal evidence of the Russia-born Ivan Ivanitz Chabert (1792-1859), who supposedly entered an iron oven in Paris (Times telescope, 1831) with a few trays of meat, and stayed there unscathed for 5 minutes at a temperature of 380 degrees F (193 oC). When he re-emerged from the oven the meat was cooked, while the only noticeable effect on him were profuse sweating and a dramatic increase in heart rate (from 98 to 168 bpm).

This experiment shows the remarkable capacity of the human body to reduce overheating by means of evaporative cooling (sweating). However, the effects of steam are far more devastating to the human body than dry heat. Not only does it carry much more energy than hot air, it also prevents perspiration by saturating the air with water vapor.

References
- American Burn Association Educator's Guide
- Times Telescope, 1831


Why are human feces riskier than urine with regard to spreading diseases? I came across a Wikipedia article titled "Urine diversion". It explains different ways, including using special diversionary toilets, and treating the collection of human feces and urine separately as opposed to just flushing everything down the toilet, maybe something most people in well developed countries are used to. These sentences caught my attention:

Separate treatment of the two types of waste is justified since urine is nearly sterile and low in pathogens, provided an individual is healthy.[3] This means that urine can be readily utilized as a fertilizer or discharged with less risk to community.[4]
Human feces, on the other hand are high in pathogens, including up to 120 viruses and need to be treated well before it can be safely used in agriculture.
Urine diversion - principle (Wikipedia article)

So urine is nearly sterile and low in pathogens, however it qualifies that statement with the condition that the person producing that urine is healthy. I assume when they compare this to feces they do so with all other things being equal (ie., provided an individual is healthy).
It claims that there may be up to 120 viruses in feces, as opposed to urine which it characterizes as being relatively low risk. I assume that feces are riskier in transmission of diseases because it naturally contains more pathogens, but it doesn't explain what these viruses or pathogens are. I was wondering if someone could elaborate on the risks of feces over urine, and also, notwithstanding that it says that urine is relatively low-risk, what health issues can be caused by reusing urine.

Many gut microbes are pathobionts or opportunistic pathogens — these organisms can be part of the normal gut flora (microbiome) of healthy individuals, but under the right (for them) circumstances (e.g.s: when introduced into the bloodstream by a cut, or due to immunosuppression) they can cause disease.1,2

As for urine, there are many fewer microbes present (>400 species and 109 to 1011 bacteria/g of contents in the large intestine vs. <100 species and on the order of 105 bacteria per mL of urine).3-5  Note, however, it is still far from sterile!5-7

You might also find the following of interest:
https://www.who.int/water_sanitation_health/dwq/iwachap5.pdf

References:
1: Hornef, M. (2015). Pathogens, commensal symbionts, and pathobionts: discovery and functional effects on the host. ILAR journal, 56(2), 159-162.

2: Price, L. B., Hungate, B. A., Koch, B. J., Davis, G. S., & Liu, C. M. (2017). Colonizing opportunistic pathogens (COPs): the beasts in all of us. PLoS pathogens, 13(8), e1006369.

3: Davis, C. P. (1996). Normal flora. In Medical Microbiology. 4th edition. University of Texas Medical Branch at Galveston.

4: Barnett, B. J., & Stephens, D. S. (1997). Urinary tract infection: an overview. The American journal of the medical sciences, 314(4), 245-249.

5: Lewis, D. A., Brown, R., Williams, J., White, P., Jacobson, S. K., Marchesi, J., & Drake, M. J. (2013). The human urinary microbiome; bacterial DNA in voided urine of asymptomatic adults. Frontiers in cellular and infection microbiology, 3, 41.

6: Hilt, E. E., McKinley, K., Pearce, M. M., Rosenfeld, A. B., Zilliox, M. J., Mueller, E. R., ... & Schreckenberger, P. C. (2014). Urine is not sterile: use of enhanced urine culture techniques to detect resident bacterial flora in the adult female bladder. Journal of clinical microbiology, 52(3), 871-876.

7: Whiteside, S. A., Razvi, H., Dave, S., Reid, G., & Burton, J. P. (2015). The microbiome of the urinary tract—a role beyond infection. Nature Reviews Urology, 12(2), 81.


How do non-human apes cool their body temperature in warm climates? Humans have sweat and we all know how that provides cooling along with the remaining vellus hairs acting as wicks.

So how do hominids, our close relatives (some of which are more massive than humans) regular their body temperature, given that they are completely covered in fur like typical mammals?

I’m specifically interested in how a human who had an atavism※ that’s like generalized hypertrichosis terminalis would fare and get by.

※ in my story, it’s due to a repaired chromosome 17 that didn’t turn out the way that was intended.

First they have less cooling needs than us, shady forests are cooler than sunlit savannah, they also are not as active, especially the larger ones, at least not in the hottest parts of the day. keep in mind chimps are much smaller than us so they have an easier time cooling off. They cool themselves the same way we do, by sweating, its just not s effective as it is in us, but then it doesn't need to be. .


Do bigger animals have more testosterone? I was wondering the other day outside, thinking about varieties of species and their aggression levels, etc.

I had a theory that since these highly aggressive, powerful, and viciously advantageous animals in fights with humans, such as lions, tigers, rhinos, elephants, etc. are so strong and aggressive in the wild compared to us humans and smaller animals, do they have more testosterone?

Like, in animals, does testosterone correlate with size, possibly explaining why a lion is heaps more slender, muscular, fit, agile, aggressive, and vicious than a typical human?

Basically, does the bigger size and aggression in wild animals correlate to their higher testosterone, such as larger size equal to larger testosterone in comparison to other smalelr or less aggressive animals, or is it different?

It is not that easy. Off the top of my hat I can try to pinpoint a few fallacies in your theory:

Scaling of the body happens in three dimensions and thus to a power law. As such allometric equations have to be applied when dealing with the total amount of substance in the body. Different species have different correlation factors of body mass to total substance amounts.
Angiogenesis (the formation of blood vessels) is affected by Allometry (the relationship of body size to the 'biology')
Diffusion factors are similar in closely related species. This means in order for messenger molecules with a short acting range/half-life (e.g. nitric oxide), the total amount in the animal may indeed be likely higher and the messenger in question may be produced on-site (e.g. endothelium). Testosterone as the name indicates, is however produced in the testes and released in the blood stream.
The sensitivity and regulation of the receptors themselves (which react to a substance) can be different in different species. It may also be weakly responsive to begin with at the genetic and ultimately the folded protein level. A good example of that is the melanocortin receptor MC1R in a gross section of the white human population. Even if you were to produce the same level of melanotonin, the melanocytes wouldn't respond to the same degree as in Humans with darker skin.
Messenger molecules can be re-purposed and take over or be given additional roles and functions over evolutionary-time. Again, Melatonin may serve here as an example.
After all a given messenger is already evolutionary tested, especially in regard to cross-talk, Haptens, and other adverse molecular interactions in the animal species and comes with the full set of molecular machinery adapted to it, - a "luxury" a new messenger doesn't have. As such pressure usually exists to re-use given messengers.

  Like, in animals, does testosterone correlate with size, possibly
  explaining why a lion is heaps more slender, muscular, fit, agile,
  aggressive, and vicious than a typical human?

I doubt you will find typical humans or typical human populations from a biological perspective. Your view may be skewed by economic factors.

  Basically, does the bigger size and aggression in wild animals
  correlate to their higher testosterone, such as larger size equal to
  larger testosterone in comparison to other smalelr or less aggressive
  animals, or is it different?

To a large degree the behavior is determined by the niche the animal is in. As to which concerted molecular messengers make the "evolutionary cut", by best influencing an animals behavior in its given "time-domain" is another question. As the question is complex much is still researched.

Take this counter-example: A mouse and an elephant have very different masses resulting in different inertial forces. Fleeing from an enemy will take different strategies in regard to the central nervous system as well as the body, and demand different levels of messenger molecules at different places, at different times for both animals. Yet the substances in question will likely be similar.


Energy transformations in chemiosmosis and ATP synthase In terms of energy, how does chemiosmosis drive ATP synthase? How does electrical energy turn into mechanical energy and then into chemical energy? Would the movement of $H^+$ be considered passive diffusion or secondary active transport?

ATP synthase is basically a pump operating in reverse: while transporters such as the Na/K ATPase uses ATP to push ions against a chemical gradient (resulting in potential energy), ATP synthase uses the H+ gradient across the mitochondrial membrane to make new ATP.

A pump uses ATP (chemical energy) to achieve one-way transport by inducing conformation changes in the protein (mechanical energy), which either opens the channel only when ions come in contact with the low-concentration side, or closes the channel when ions come in contact with the high-concentraton side. In either case, the result is that ions can pass only in one direction, and a concentration gradient will build up (potenial energy).

In the case of ATP synthase, the process is simply reversed: an H+ ion entering the channel from the high-concentration side (the intermembrane space) causes a conformation change that opens the channel, and is coupled to the synthesis of ATP. Precisely how the ATP synthase accomplishes this on the molecular level is a complicated topic, and still an active area of research. See this article for an overview.

Note that ATP synthase is clearly reversible: if the mitochondrial H+ gradient collapses, for example if respiration is blocked, then ATP synthase can run in reverse, consuming ATP to pump H+ out of the matrix. There is even an ATPase inhibiting protein dedicated to blocking the ATP synthase in this situation, to prevent the depletion of ATP. Also, the lysosomal ATPase is structurally similar, but operates in the reverse direction to create low pH in lysosomes.


Identification of a young tree in the alpine region of southeastern France I have a young tree growing in my backyard (in Grenoble, France), and I was not able to identify it, despite searching with many identifier tools.

It's 2 years old now and must be 2 meters high. It has a grey trunk, its leaves grow on brown/red stems that protrude irregularly from the trunk.

The leaves are pointed, gently serrated, covered with a light fuzz, with a distinct lattice of veins.

Could you help me identify it?

The tree grew its first fruits this spring, so that I could definitely identify it: it is a goat willow! This is consistent with its spontaneous sprouting in my yard, as it is a common tree in my region.


What is this arthropod? It is approx 2cm long and was discovered in house after a bush walk & BBQ at Riverside Garden Reserve in Mandurah, Western Australia. It has 6 legs.

Thanks! :-)

This is a Sawfly larvae, and is most likely of the species Perga dorsalis or Perga affinis. The larvae you depict appears to be quite young, but as it matures it'll gain a much darker color, and, depending on whether it's P. affinis vs. P. dorsalis, the abdomen will grow fairly long.

I imagine the developmental progress of your larvae is somewhere between the two following larvae.

A few more images:


Is it possible to correctly identify presence of Y chromosome with external physical test only? I asked a question related to the third sex, and I came to know that its always possible to categorize a human to male or female with presence of Y chromosome.

Now, I have another question. Is there a way to say if someone has Y chromosome only with external physical test? More specifically:

Someone has a penis, can we say that s/he will have Y chromosome?
Some one has Y chromosome, can we say that he will have a penis?

No, an external physical examination would be inconclusive. The reason is the TDF
gene. To be more specific, if a person has XY and the gene is not active then the subject
would have a female appearance. Also we cannot conclude that a person has a Y
chromosome even if it has penis because that gene could be transferred on the X chromosome.

Here is a reference: http://en.wikipedia.org/wiki/Testis_determining_factor.


Are new blood vessels able to deliver oxygen to hypoxic tissue before the establishment of blood circulation? Sprouting angiogenesis - the growth of new blood vessels from a preexisting vasculature - can be triggered by cells in hypoxia in order to re-stablish the oxygen and nutrients supply to the tissue.
The formation of a new vessel is a multistep process starting by the activation of the Endothelial Tip Cell (ETC) phenotype that migrates through the tissue guided by chemotactic cues, such as the Vascular Endothelial Growth Factor (VEGF) concentration. The cells near to the ETC gain the Endothelial Stalk Cell (ESC) phenotype and start to follow the ETC and proliferate in order to form the new branch.
The Endothelial Cells (ECs) re-arrange in the new branch, forming the lumen. Thus Smooth Muscle Cells (SMC) and pericytes are recruited to maturate the new vessel. Consequently the vessel is filled with blood from the main branch.
My question is: when does the new vessel starts to deliver oxygen to the surrounding tissue? Can it occurs before the anastomosis formation where the circulation is stablished, even for a short period of time due to the plasma that filled the new branch? (See the following picture where they represent the cells near to the new branch with a high concentration of oxygen even without anastomosis/blood circulation)

Image Source: Peter Fraisl et al. Regulation of Angiogenesis by Oxygen and Metabolism. DOI 10.1016/j.devcel.2009.01.003

It is true that for long-term oxygen delivery the new vasculature must establish circulation. However, when the new vessels are starting to maturate they are filled with blood coming from the main vessel, which carries oxygen.

Since the endothelial cells of the new vessel derive most of their energy anaerobically through glycolysis (Quintero et al, 2006; Dagher et al, 2001; Krutzfeldt et al, 1990; Peter Fraisl et al, 2009), a metabolic mechanism  activated in critical microenvironments like hypoxic tissue (DeBerardinis et al, 2008), the oxygen in the lumen is free to be immediately delivered to the tissue.

Thus it is fair to assume that  at least for a short period of time the cells near to the new vessel are being oxygenated.

References:

Quintero, M., Colombo, S.L., Godfrey, A., and Moncada, S. (2006). Mitochondria as signaling organelles in the vascular endothelium. Proc. Natl. Acad. Sci. USA 103, 5379–5384.

Dagher, Z., Ruderman, N., Tornheim, K., and Ido, Y. (2001). Acute regulation of fatty acid oxidation and amp-activated protein kinase in human umbilical vein endothelial cells. Circ. Res. 88, 1276–1282.

Krutzfeldt, A., Spahr, R., Mertens, S., Siegmund, B., and Piper, H.M. (1990). Metabolism of exogenous substrates by coronary endothelial cells in culture. J. Mol. Cell. Cardiol. 22, 1393–1404.

DeBerardinis, R.J., Lum, J.J., Hatzivassiliou, G., and Thompson, C.B. (2008). The biology of cancer: metabolic reprogramming fuels cell growth and prolif- eration. Cell Metab. 7, 11–20.


What is the name of this small black insect? Location: South India

Size: 1-2 cm

Precise identification of insects requires use of microscopic lens and an immobile specimen.
I have neither of those here; also, my reference material is focused on North American species.  Nevertheless, with what I have available to me, it is most likely a member of order Hemiptera, family Alydidae:

Alydidae are mostly active during the day, so your specimen was probably resting and saw a nearby artificial source of light.  They use their piercing mouth-parts to suck the juices out of herbaceous seeds.  Unless yours was discovered in a heavily urbanized area, I would venture that this one feeds off nearby grasses and trees.


Cells scatter light in solution. Do they refract? For part of an exam question, "why do cells grown in liquid cause turbidity?" I answered that cells refract light. The correct answer was that cells scatter light. Isn't the light scattering ultimately caused by lots of refraction events? Is the grader correct to take off full points?

You're both correct. Turbidity is causes by excessive light reflection events from all of the particles, and is actually how it is measured in a sample. As to cells, they both both reflect light that is incident on a membrane, but they will also refract light that does manage to enter the cell because of the difference in refractive indices. A portion of the light will also be absorbed by the cells.


Software for counting fly eggs Is there any software that could be used to rapidly (quicker than humans) count the eggs laid by Drosophila on a substrate by placing them under a microscope with a camera attached? The eggs are laid on the surface of the food and have a fairly uniform shape - we regularly have to count the number of eggs in >200 vials, which can take a over 4-5 hours per 100 vials.

I wonder if there is some software which could be linked to a microscope with mounted camera connected to a computer, the software could then count from the live feed or a quickly taken photograph. It wouldn't have to be perfect but as long as it could be as accurate as humans (which is surprisingly low from our experience - we count up to a set amount then remove the rest, e.g. if I was aiming to leave 180 eggs I would be quite happy to get anywhere between 150-200 eggs). It would be using a stereo-microscope (aka dissection microscope) with a mounted camera.

The surface of the food with eggs looks something like this ->

purely for illustration, taken with iPhone - real images on microscope would be able to get closer (food surface could fill screen) and have more uniform lighting to give better contrast (this was taken to show in a talk so I wanted pretty rather than practical). The bottom two quarters of the food in this picture would be similar to the real result, just much closer...

ImageJ is a multi platform piece of software that has a cell counter module that might be of some use, and hey its free! Its easy to use and so ludicrously crude that it is very  versatile. I remember using it in undergrad to count cells under the microscope automatically after a few image contrast tweaks. I don't see why this couldn't be reapplied to count eggs.

Download: http://rsb.info.nih.gov/ij/download.html

Cell Counter documentation: http://imagej.nih.gov/ij/plugins/cell-counter.html


Why does cell waste energy in meiosis, between meiosis 1&2 As far as I have learnt about meiosis I have read that anaphase 1 is followed by telophase 2 where chromosomes change back to reticulum.but in the very next stage i.e prophase 2 they again start condensing to form chromosomes again.i have also heard that is most species, cytokinesis is simultaneous and not successive( which means that the 4 cells undergo cytokinesis together)..

If this is the case, then what is the need for a cell to waste its energy on making chromosomes back to reticulum at the end of meiosis 1, and again condense it in the very next step
Can't it directly continue on to prophase 2!

It sounds ridiculous,...but I beleive there must be a reason behind it, because the last thing a cell would want to do is to waste its energy in such things.

Thanks in advance

There is a phase between meiosis 1 and meiosis 2 called Interkinesis.
During this phase centrosomes or centriole pairs undergo replication in animal cells, which is important for bringing true haploidy.

Refer - meiosis-ii


Why is it common to grow microalgae in bottles or canisters? I came across many people growing Chlorella or Spirulina types of microalgae, inside bottles or canisters and that that they also bubbled CO2 into these containers artificially.

This question is basically comprised of two questions that will clarify what I am trying to understand:

Why do they use bottles or any type of bottle-like canister?
Why do they insert CO2 to the bottle artificially in cases the bottle doesn't have a cork ? If the claim is that enough CO2 won't be able to get through the bottle neck than why growing the algae inside a bottle\aquarium in the first place, and not in a tub with full exposure to the air?

Really, the credit for this answer goes to another 'Homo sapien' from the Spirulina manual link in his comment.

Why do they use bottles or any type of bottle-like canister?

Likely, because they are clear. From the manual:

Spirulina needs sunlight so it is preferable that the container in which it is grown be transparent.

Plastic bottles are cheap, clear, and easy to come by, so they make the ideal material.

Why do they insert CO2 to the bottle artificially in cases the bottle doesn't have a cork?

I don't think the tubes in your picture are CO2 tubes. Rather, I'd bet they are bubble tubes (like in aquariums). These are necessary to keep the growth medium moving constantly. Again, from the manual:

Spirulina tends to gather at the top of the growing culture, where sunlight exposure is maximal. Due to this, Spirulina that cannot reach the top will not multiply and will ultimately die.
In order to maximize Spirulina exposure to sunlight, the water in which it is grown must be stirred
...
Another option is a pump, the simplest kind used for aquariums.
(emphasis mine)


What is the name of a gene whose effect comes from the dominant and recessive alleles The p53 gene has a mutation which results in the cell having a cancerous phenotype.  This effect is dominant as one allele being mutated is enough to result in this effect.  However when both alleles are mutated then the phenotype changes / become more pronounced so it's not a simple dominant mutation.  What is this kind of mutation called?  Is it hemi-dominant or something?

Allelic effects

  What is the name of a gene whose effect [..]

The concept of dominance make sense for alleles. The gene has no name to refer to the allelic interaction.

Patterns of dominance

Patterns of dominance and recessivity are typically used for quantitative traits. Can you somehow put a a number into the phenotypes of all three genotypes? Making some kind of association table would help even if the phenotypes are not quantitative. See below.

Attempt to answer the question

It looks like you are considering a case where the phenotypes

aa -> healthy
Aa -> sick
AA -> very sick

Under this model, it is hard to tell whether you would consider sick half way through healthy and very sick or not. It feels like you want to consider sick as being closer to very sick than to healthy so if we put numbers we could get something like

aa -> 1 (healthy)
Aa -> 2 (sick)
AA -> 2.5 (very sick)

As such, A shows partial dominance (aka incomplete dominance or semi-dominance).


How many mRNA strands are in a single dose of the COVID-19 vaccines? I realize there are several different mRNA vaccines. I would be happy to know the ballpark figure for any of them.
As a follow-up, is it known about what percentage of injected mRNA strands are successfully translated into proteins?

Well, we can work this out empirically.
We know that the average molecular mass of a base of a RNA is 339.5 g/mol.
However to calculate the approximate molecular mass of a single-stranded RNA molecule, you multiply the number of bases by 320.5 and add 159 (same source). This means for a single-stranded molecule of 1000 bases, it will have a molecular mass of 320,659 g/mol. The Pfizer vaccine has a sequence of 4284 bases.
4.284 * 320659 = 1373703.156 g/mol.
The Pfizer vaccine has 30 ug/shot (PDF, see page 27 under "Description") or 30x10-6 grams
number of moles = mass/molar mass.
n= 30 x 10-6 grams/137370.156 g/mol
n= 2.18 x 10-11
number of molecules = n (above) x Avogadro constant (6.022 x 1023)
Number of molecules = 1.3 x 1013
Now this is a ball-park figure because the RNA isn't actually RNA it has had some special additions to the uridines that make it less susceptible to immune regulation, replacing the uridines with  1-methyl-3'-pseudouridylyl (see Wikipedia link to structure above). However, I believe this alteration to the mass, will not alter the number substantially; it'll still be in the 10 trillion range.


What is the percentage of people living in England in 1500 AD whose lineage is still alive? This sounds a bit random, but it stems from a lecture in statistical genetics which I attended a while ago. We were shown a population lineage graph from which it was clear that most lineages eventually go extinct. The further you go in to the past, the smaller the percentage of the individuals living at that time whose lineage is still alive today. This probably makes sense in the context of evolution.

So I am curious about the speed of that process in human populations. I am aware that this speed, quite possibly, varies over time and geographically. Are there any estimates about that speed in European populations? To be more specific, I would take England at 1500 AD as a starting point. What is the percentage of the people living in England at that time whose lineage still exists?

EDIT

I am very grateful to @Richard Erickson for the reference to the Galton-Watson process. It is really fascinating that +70% of all known Chinese family names have become extinct. It would be really interesting to know more about the speed of the process of extinction - in particular in European populations.

Here is a partial answer that also expands upon a previous comment:

Between 1900 and 2013, 20,000 English surnames had gone extinct according to the Daily Mail's summary of research on surnames. This is a partial answer because it only tells us how many names are gone and is from the wrong time period.

To completely answer your question, I suspect you need to conduct primary research. I would suggest the following broad steps:

Gain access to the data needed to answer your question. British Census records only appear to available from 1841 to 1911 so you most likely need to sample old names (e.g., graveyard based studies).
Use the data to parameterize a statistical model such as a Galton-Watson branching process models.

I doubt the answer to your question can quickly be found unless somebody has already done it. Although, if done well, it could probably be a masters' thesis. However, before embarking on such a project, I would suggest doing a good literature search to make sure it has not been done before.

Note: If anybody finds the answer, please feel free to take content from my answer and include it yours.


Were nicotinic acid/amide or flavin nucleotides ever part of primary RNA sequence? NADH and FADH2 redox reactions are built deep into our biochemistry.  For example, pyridine nucleotides are involved in >500 enzymatic reactions.  When we look at the structure of deamido-NAD+, it is strikingly symmetrical:

The same is true for FAD, a two-headed monster of a nucleotide

in which we can recognize the two heads, belonging to xanthosine (a purine nucleotide) and alpha-ribazole 5'-phosphate, the latter of which is produced from flavin in the course of B12 synthesis:

The phosphoribose can be swapped between alpha-ribazole 5'-phosphate and nicotinic acid.
Given such enticing hints, it is tempting to speculate that one or more of these nucleotides might be so fundamental because it was actually part of some ancient RNA primary sequence, giving RNA containing this code the power to do serious catalytic chemistry with acid and amide groups, powerful options for redox and methyltransferase activity and more.
I know this is a hard question to answer, but is there any evidence that weighs in, or could weigh in, on whether this happened?  For example, if these functional groups are too unstable in regard to intramolecular reactions on a strand of RNA, or if there is no way to envision a potential base pair that would be reliably distinguished from the known base pairings.  Or if RNA editing sometimes creates them in some organism.  Or if aptamers made from these subunits show improved catalytic possibilities.  (That doesn't seem so easy to look up, since many papers look for aptamers binding these as prosthetic groups, rather than incorporating them)  It would be useful simply to have a notion of what evidence for such an idea would look like.

There are naturally occurring RNAs that contain a NAD "cap" at 5' end (Bird et al., 2018; Wang et al., 2019).
I could not find any papers on natural RNAs with nicotinamide or flavin as a regular base. However, a study with synthetic oligonucleotides reports that NAD containing oligonucleotides can form stable triplex structures (Göckel & Richert,2015).
There are two studies from 1970s that explore base pairing potential of flavin. Raszka and Kaplan (1974) report that FMN can hydrogen-bond with AMP but not other nucleotides. They further explore this finding to show that FAD can form intramolecular hydrogen bonds. Scarbrough et al (1976) show that flavin containing nucleotides can form stacking interaction in addition to base pairing kind of hydrogen bonds. (Compare D-X bonds)

More recently, Julius et al. (2020) have shown that NAD and FAD can act as co-factors of primase to initiate DNA replication.
Overall, it seems plausible for FAD and NAD to be a part of nucleotide chains (more plausibly for RNA as they are already known to contain a variety of unusual bases). However, I cannot comment if they really did play such a role during cellular evolution.


What is the strand specificity of a reference genome? It's a simple question but I've come across many people who have this question, is the reference genome Positive of Negative strand? Indeed, I've had heated arguments over the same issue.

So here's to putting all those questions to rest.

Is the reference genome a positive or a negative strand? And why so.

First of all, the reference genome strand specificity is referred to as sense (positive strand) or antisense (negative strand). Now let's consider to sequencing data or FASTQ files. When we align reads, a resulting SAM or BAM file has a column specifying strand information, we usually see a + or - strand.

For more background on the strand names here's a somewhat incorrect note from the wiki article under sense (molecular biology)

  The strand names actually depend on which direction you are writing
  the sequence that contains the information for proteins (the "sense"
  information), not on which strand is on the top or bottom (that is
  arbitrary). The only real biological information that is important for
  labeling strands is the location of the 5' phosphate group and the 3'
  hydroxyl group because these ends determine the direction of
  transcription and translation. A sequence 5' CGCTAT 3' is equivalent
  to a sequence written 3' TATCGC 5' as long as the 5' and 3' ends are
  noted. If the ends are not labeled, convention is to assume that the
  sequence is written in the 5' to 3' direction. Watson strand refers to
  5' to 3' top strand (5' → 3'), whereas Crick strand refers to 5' to 3'
  bottom strand (3' ← 5').[4] Both Watson and Crick strands can be
  either sense or antisense strands depending on the gene whose
  sequences are displayed in the genome sequence database. For example,
  YEL021W, an alias of URA3 gene used in NCBI database, defines that
  this gene is located on the 21st open reading frame (ORF) from the
  centromere of the left arm (L) of Yeast (Y) chromosome number V (E),
  and that the expression coding strand is Watson strand (W). YKL074C
  defines the 74th ORF to the left of the centromere of chromosome XI
  and denotes coding strand from the Crick strand (C). Another confusing
  term referring to "Plus" and "Minus" strand is also widely used.
  Whatever the strand is a sense (positive) or antisense (negative), the
  default query sequence in NCBI BLAST alignment is "Plus" strand.

Although it is correct when it says

Watson = Sense = Plus Strands
Crick = Antisense = Negative Strands

A satisfactory answer/convention is provided/suggested in this PMC article covering the terminology of the Watson and Crick strands

Below are some interesting excerpts from the publication:

  The earliest reference that we could find to the "Watson strand" and
  the "Crick strand" is somewhat tongue-in-cheek and comes from a pair
  of papers in 1967 by Wacław Szybalski and colleagues. They
  bound the two DNA strands of phage λ to the synthetic polynucleotide,
  poly(IG), which has an affinity to cytosine-rich regions. They then
  separated the two strands by density, which turned out to be
  determined by the amount of bound poly(IG). In a cesium-chloride
  density gradient, the strand with more bound poly(IG) was denser and
  heavier than its complement. Because the "dense" strand was
  cytosine-rich, Szybalski and colleagues called it the "C strand."
  Logically, thus, the complementary strand, which was guanine rich,
  should have been the "G strand." Instead, it was christened the "W
  strand."

So the authors note that the Strands started on their terminology journey with their roles reversed with the Crick strand specifying the heavy strand (with IG) and the Watson strand specifying the lighter one.

They also note that the present day terminology of the WC model is non-arbitratrily based on the current horizontal drawing convention where one strand is placed on top and the other at the bottom.

They also suggested a typical convention which is followed today. If it was due to their suggestion is a different matter altogether, but beyond the scope of this question.

  Given the amount of effort already spent on standardizing such
  databases, and their influence on other disciplines, we feel that the
  genomic definition of Watson and Crick strands has the most mass
  behind it. Specifically, we find that the unambiguous usage of the
  Saccharomyces Genome Database to be the most useful. Under the first
  part of our proposal, the centromere is a reference point that divides
  a chromosome into two arms of unequal length. The chromosome is
  oriented so that shorter arm is on the left and the longer arm on the
  right. Furthermore, the top strand has its 5'-end at the left
  (short-arm) telomere and its 3'-end at the right (long-arm) telomere.
  This strand is the Watson strand. Similarly, the bottom strand has its
  5'-end at the right telomere and its 3' at the left telomere and is
  the Crick strand. We further propose that "top", "forward", and "plus"
  be used as synonyms for the Watson strand and "bottom", "reverse", and
  "minus" for the Crick strand.

and later on...

  If it is ultimately impossible to distinguish Watson and Crick strands
  using biological properties, then we propose that Watson should refer
  to the stand arbitrarily used as a reference in a database (i.e. the
  "plus" stand) and the Crick strand should refer to its complement

But let's come back to sequencing data or FASTQ files. When we align reads, a resulting SAM or BAM file has a column specifying strand information, we usually see a + or - strand.

Which suggests that the product originated from either the Watson (positive) or the Crick (negative) strands. For example, a read is actually the reverse complement of the product, but since you do a PCR step during you library preparation, that particular bit of information is lost and therefore the downstream analysis protocols tend to consider the entire loci any single product aligns to.

The difference between strand specific sequencing and non-strand specific sequencing are covered here.


Why are sugars like mannose expressed on the outside of eukaryotic cell membranes? Bacteria are able to adhere to sugars (e.g. mannose) on the exterior of eukaryotic cells, leading to infection and disease. Why have eukaryotes not evolved so as to dispense with sugars on their cell surfaces? I assume there must be a positive aspect to having these sugars. What is it?

The “positive aspect” or reasons why the surfaces of eukaryotic cells contain sugars are various. Major ones appear to be cell adhesion and cell–cell interaction.

In fact, the variation of the sugars in certain cell-surface antigens (blood group antigens and histocompatibility antigens may represent the kind of evolutionary response to infectious agents mentioned in the question. The situation with the major histocompatibility antigen is complex, but variation in blood group antigens is thought to be a response to the malaria parasite — reviewed by Cooling.


Is RFLP markers are dominant or codominant? there are several markers used in molecular biology,. marker analysis. RFLP is most commonly used marker which having several advantages. is RFLP marker co-dominant or dominant ?

I suspect that you "know" the answer to your question and are just seeking confirmation.  First, let's get some terms straight.  Dominant and recessive are used in reference to phenotypes, so one needs to have a transmissible phenotype to score among the offspring of our genetic cross.  Phenotype are usually due to (or caused by) mutations.  In a diploid species, if you can only score the phenotype in offspring that are homozygous for that mutation (denoted m/m) then the phenotype is recessive to the wild-type allele of the gene.  In other words, a heterozygous offspring, (m/+) will have a wild-type phenotype.  On the other hand, if a heterozygous (m/+) offspring has a mutant phenotype then we would say that the mutation is dominant over the wild-type allele.
Your question is asking about a test cross of a +/+ individual to a homozygous mutant (m/m), however, in this case the RFLP does not typically have any mutant phenotype, it is just a convenient genetic marker for the purposes of this discussion.  All of the cross-progeny will be (m/+), and if you go to the trouble of preparing genomic DNA from those F1 cross-progeny and perform a Southern blot with a probe that detects the RFLP, then you will be able to demonstrate this.  Typically, unless the RFLP-bearing chromosome is homozygous, the probe will also detect the wild-type restriction fragment, so there will be two bands on the film, one for each allele.  Since you can detect the genetic marker in the F1 generation, the RFLP is dominant, but since you can also "score" the wild-type fragment, the two fragments are indeed co-dominant (after you do the Southern blot).


how to change format of ent file into PDB format? whenever i download files from Protein Data Bank(PDB),they are downloaded in .ent format. How can i convert them into PDB format??

Just rename the file to .pdb - in other words, .ent files are PDB format:

  The archival PDB files will be distributed with the reserved conventional names, in the form pdbentry_id.ent, where entry_id is a PDB 4-letter code, e.g. pdb1abc.ent, for PDB format entries; rentry_idsf.ent, e.g. r1abcsf.ent, for X-ray experimental data; entry_id.mr, e.g. 1abc.mr, for NMR experimental/constraints; entry_id.cif, e.g. 1abc.cif, for mmCIF format entries; and entry_id.xml, e.g. 1abc.xml, for canonical XML format entries.

From http://www.wwpdb.org/about/faq


Understanding the configuration of channels in the 10-20 EEG montage In a typical 10-20 system there are 21 electrodes placed on the scalp. However, this does not mean there are 21 distinct "channels" or voltage sources.

I've heard that some of these electrodes are "reference" electrodes, while others are "active" electrodes. Sometimes electrodes can setup in a "bipolar" or "differential" fashion.

How many actual channels (that is, distinct input sources) will there be?
What's the difference between reference and active channels?
How does this bipolar/differential setup work differently than a "normal" setup?

Update

Here is my updated understanding based on @Christiaan's answer below:

time (t)    F2v    F3v        F2-F3         F3-F2
=================================================
1           2      1            1             -1
2           3      6           -3              3
3           5      3            2             -2

So, if my understanding is correct, then at time t=1, the voltage of the F2 electrode might be, say, 2 (units; volts, microvolts, whatever), and the voltage of the F3 electrode might be 1. If F2 is the arbitrary active electrode, then the potential difference between F2 and F3 is 2 - 1 = 1. But, if F3 was the active electrode (and F2 the reference), then the potential difference would be the inverse (that is, 1 - 2 = -1).

EEG 10-20 system. source: Wikipedia

  How many actual channels are there?

21 in the figure, i.e., the number of active electrodes.

  What's the difference between reference and active channels?

The active electrode is the electrode under investigation. Basically this is arbitrary. Consider electrode F2 and F3. When you measure the potential difference between them and F2 is the active electrode and F3 the reference, the signal will be exactly identical as in the reverse situation, only the polarity (the sign of the voltage) is reversed. That's all there is to it. The active electrode is typically the electrode linked to the location you wish to record. The reference is elsewhere.

  What is the difference between a bipolar and a normal
  (unipolar/monopolar) setup?

In a regular unipolar setup, each of the electrodes is either measured against a distant reference (e.g. the earlobe) or against the aggregate of all the electrodes. Either way, the reference is distant. This means that the signal will be high, but artifacts will be large too, because artifacts from both the electrodes will be merged into the signal. Consider a reference in the neck; neck musculature activity will add artifacts (EMG) to the signal. Likewise, a reference in the chest area will add ECG to your signal.

In a bipolar setup adjacent electrodes are recorded. This means that large artifacts occurring through eye blinks or whatsoever are recorded by both electrodes and as their polarity is opposite, the large artifact responses will be cancelled and the signal is much cleaner. However, the downside of the coin is that commonalities in the signal, which will be larger when the electrodes are closer will also be subtracted and hence your signal will decrease accordingly.

I have used bipolar setups when measuring eERGs and eCAPs, i.e., electrically evoked activity. The electrical stimuli generated large artifacts and actually so large that their amplitudes were approximately similar between active and reference electrodes. The reduction in artifact by far outweighed the loss in signal. It depends on the situation you are at.

An authoritative work in this field is listed below.

Reference
- Pivik et al., Psychophysiol (1993); 30: 547-58


In SDS-PAGE, why are the sample buffer and running buffer different in concentration? I'm currently trying to understand the concepts of SDS-PAGE. Every protocol I've read says to use 5X for sample buffer and 1X for running buffer, but I don't really understand for what reason? What are the negative consequences if I don't?

The reason why the sample buffer is more concentrated (typically 2x or 5x depending on your protein concentration) is its dilution when you mix it with the sample. You mix 4 parts of your sample and 1 part of 5x sample buffer, so that the final concentration of this buffer is 1x.

If you use 1x sample buffer together with the same amount of sample, the final concentration of the sample buffer would be 0.5x, if you do the same 4+1, it would only be 0.2x in the end.

This dilution is too strong and the buffer will not fulfill its purpose, the complete denaturation of all proteins in the sample, so that their distribution on the gel depends only on their molecular weight and not also on their tertiary/quarternary structure.

Using a 5x buffer allows you to add a lot of sample when diluting it to 1x without changing the concentration of the protein too much. It also allows to load larger amounts of sample compared to 2x buffer, interesting when the protein concentration is low.

The running buffer is not further diluted by adding sample, so that it has to be diluted to 1x in order to have the right ionic strength. Since these are often used in big quantities in the lab, it is more convenient to prepare concentrated buffers in smaller quantities and dilute them before usage than preparing big tanks of buffer.


Why do viruses last longer on hard surfaces, like glass, than porous ones All the papers I've read about survival of viruses outside of the host refer to longer survival times on hard surfaces without explaination. As a layman I would have thought that a virus (who's main environmental dangers are things like UV would survive better on rough or porous surfaces, where they are protected and can hide from the light and weather, not the other way round.

For example, one paper I read recently (although it's an old paper) showed Ebola virus surviving for 50 days on glass but not even 2 on rough materials. Can anyone enlighten me as to what it is about glass that makes such an ideal home for viruses to extend its lifespan by over 25 times?

More scientific papers would be helpful, but the medical community does seem to think that hard surfaces in general preserve viruses better than soft ones. The Mayo Clinic and UK NHS web sites both make that claim, without explanation.

The paper you cited is not about survival on smooth versus rough surfaces; it does, however, mention discrepancies in survival on metallic versus non-metallic surfaces.

From the article you cite (Piercy et al., 2010):

  The lack of recovery of virus from metal substrates may be attributed to several factors; previous work (Sommer et al. 1999; Pawar et al. 2005) has shown the binding of micro-organisms to metal surfaces to be poor because of the high surface energy, high electronegative and hydrophilic properties of metal. Other research has also shown that positively charged metallic ions such as copper and silver have a high bactericidal activity (Friedman and Dugan 1968; Bitton and Freihofer 1977; Slawson et al. 1990) and hence by extrapolation may also have a high virucidal activity.

In short, it isn't exactly clear why viruses in particular might have poor survival on metal, but the electrical and chemical properties of metallic surfaces, as well as biological activity of metallic ions that are constantly leached from alloys like steel are all candidates.

The claim may also not apply to all viruses, and the results are not consistent across labs. For example, Sagripanti et al. 2010 find no effect of surface material (aluminum alloy, rubber, or glass) on Ebola, Venezuelan equine encephalitis, Lassa, and Sindbis viruses.

References:

Piercy, T. J., Smither, S. J., Steward, J. A., Eastaugh, L., & Lever, M. S. (2010). The survival of filoviruses in liquids, on solid substrates and in a dynamic aerosol. Journal of applied microbiology, 109(5), 1531-1539.

Sagripanti, J. L., Rom, A. M., & Holland, L. E. (2010). Persistence in darkness of virulent alphaviruses, Ebola virus, and Lassa virus deposited on solid surfaces. Archives of virology, 155(12), 2035-2039.

Bitton, G., & Freihofer, V. (1977). Influence of extracellular polysaccharides on the toxicity of copper and cadmium toward Klebsiella aerogenes. Microbial ecology, 4(2), 119-125.

Friedman, B. A., & Dugan, P. R. (1968). Concentration and accumulation of metallic ions by the bacterium Zoogloea. Dev. Ind. Microbiol, 9, 381-388.

Pawar, D. M., Rossman, M. L., & Chen, J. (2005). Role of curli fimbriae in mediating the cells of enterohaemorrhagic Escherichia coli to attach to abiotic surfaces. Journal of applied microbiology, 99(2), 418-425.

Slawson, R. M., Lee, H., & Trevors, J. T. (1990). Bacterial interactions with silver. BioMetals, 3(3), 151-154.

Sommer, P., Martin-Rouas, C., & Mettler, E. (1999). Influence of the adherent population level on biofilm population, structure and resistance to chlorination. Food microbiology, 16(5), 503-515.


FPLC based separation of serum proteins I am working on a project which involves bio-marker discovery in neuro-developmental disorders. I have an akta-prime FPLC instrument in my lab but I do not know how to use it.

According to theory, plasma, or serum, normally contains numerous proteins of high molecular weight — which might mask the concentration of small-molecular-weight proteins that could be relevant to my study.

I am looking for the abnormal expression, if any, of such protein molecules.
But I am concerned that if I run the serum sample directly, I will get too many bands to allow me to identify my desired molecule.

Which column should be used to separate such masking proteins from serum — or, alternatively, what previous separations should be done?

Also, what reagent concentrations should I use (particularly re: binding buffer and elution buffer)?

Firstly, may I suggest 2D electrophoresis as a technique for identifying differences in protein/marker expression levels. This technique will give you greater resolution and ease of comparing your sample in basically one step/experiment. Added bonus once you determine your protein of interest you can easily excise it and send for mass spec analysis to identify your protein. In short, 2D electrophoresis separates protein based on their pI(isoelectric point) in 1st Dimension and then by their size in the 2nd Dimension. Once you identify the protein of interest you can purify it using the below steps.

Image Source

Now onto tackling the problem of separating your proteins for analysis. The general outline you can follow to separate a protein, is a follows:

Ion Exchange Column
Hydrophobic Interaction column (HIC)
Size-exclusion Chromatography (SEC)

Ion Exchange first because you elute the protein off the column using increasing concentration of salt (0 --> 1M of NaCl). Whereas you need a high concentration of salt to load the protein onto HIC column. So, the elute from setup 1 can be easily spiked with a 5M (Nh4)2SO4 and loaded onto HIC column (No dialysis). Also, Ion-exchange and HIC are high capacity columns whose resolution can be improved with shallow salt gradients.Hence my recommendation to be used first. However, if you are certain your marker/protein of interest is small you can directly use SEC column

Ion exchange column, as the name suggests, separates protein on the basis of their charge (Least ---> Most Ionic, also note charge +/- of the protein at that pH matters). The buffer to use for Ion exchange would be:

 1. Binding Buffer
    10-50mM Phosphate Buffer pH7.4 (plasma pH)
    0-25mM NaCl
 2. Elution Buffer
    10-50mM Phosphate Buffer pH7.4
    1M NaCl

Use low salt as possible to get the most efficient binding to the column.
Elution is done with a gradient of NACL from 0% of elution buffer to 100% Elution over about ~20 column Volumes.

HIC column separates the protein on the hydrophobicity of the protein. The order of the elution from the column is least --> most hydrophobic.

1. Binding buffer
    10-50mM Phosphate Buffer
    1M (Nh4)2SO4
 2. Elution Buffer
    10-50mM Phosphate Buffer
    (No Salt)

You can find a good protocol for HIC here. The elution is also done over ~20 CV with a gradient of 0-100% elution buffer, similar to Ion exchange elution.

Finally, size Exclusion or SEC separates proteins based on their sizes. The elution order is biggest ---> smallest size. The column is run in an isocratic mode i.e only one buffer used.
Buffer composition:

1. 50-100 mM Phosphate Buffer
2. 150mM NaCL

The salt in buffer is to prevent non-specific interaction of the protein with column, which might result in protein eluting later than necessary. This might result in poor resolution.

The most important thing to realize is each of the above column have their own sets of cavets. 1) Ion exchange columns are specific to the charge of protein, so you will end up having a lot of protein not binding to the column at all based on type (anion/cation) of column used. 2) Similarly, HIC column might not bind well to your protein as well. Also, adding 1M (Nh4)2SO4 will precipitate proteins which might include your protein of interest. 3) SEC columns have an upper limit on protein size they can accept, above which all the protein elute together at the void volume decreasing the resolution(However, this shouldn't be a problem in your case as you trying to separate out the lower molecular weight protein from the higher weight ones.)

Finally, all the above step work great considering you know the protein you are looking to purify(separate). Since that is not the case you will have to collect and run SDS-gel Electrophoresis of the fractions at each column purification stage. Also, run a side by side comparison of control and sample plasma to figure out which protein is being expressed differently.

In closing, I would like to point out that 2D electrophoresis gives you a gel to compare, as shown in the pic, anyway. So it might be a better solution to your problem. However, I do admit you might have to setup purification to confirm your findings.

Hope that helps, good Luck!


Remove all intestinal bacteria? I've heard that this has happened to a few people before in extenuating circumstances, but don't know where to search for more information on this.  Is there a name for this?  What are some major symptoms, and how does it affect the body?  How is it identified?

This has never happened (at least to my knowledge), since removing all bacteria from a human being would be extremely difficult and also dangerous. And there is a proven relationship between the human microbiome and stress, but "extenuating circumstances" won't remove all bacteria from your body.

However you can look in the Internet for "germ-free animals", "gnotobiotic animal" or "axenic mice". Caesarean section in aseptic conditions is practiced (a fetus in normal conditions has no germs in it) and the "newborn" is immediately put in a sterilized enviroment. There is plenty of research with this kind of animals.

I hope my answer was useful for you! Cheers!


Do bacteria with multiple flagella move faster than bacteria with a single flagella? Do bacteria with multiple flagella move faster than bacteria with a single flagella? Assuming the flagella are at the same length.

Note: in the following I'm going to focus on E. coli, which is the best understood system for bacterial flagellar movement. It would be interesting to compare uniflagellate bacteria and bacteria that have a fixed number of flagella. However, I know about E. coli and I think there is something interesting to be learned just from them, so I'll continue.

You're in luck that you asked this question when you did: a very cool paper (1) was published just last year addressing your question. They find that the number of flagella has only a small effect on bacterial motion, but to understand why requires some background (A lot of the following background comes from this paper, which summarizes it pretty well).

I would argue that E. coli movement is one of the best-understood examples of cell movement we have, with excellent, high-resolution, quantitative data being generated for over 40 years now (the king in this field is Howard Berg at Harvard, if you want to learn more just read a recent review from him).

E. coli have a variable number of flagella, from one to ten or so. They operate in a bundle, but this bundle is highly dynamic due to the way E. coli moves. Rather than proceeding continually a straight line like a car, E. coli moves by a series of starts and stops called a run-and-tumble: the bacteria moves in a straight line for a short amount of time (a run), and then it stops in its tracks and changes direction (a tumble) before starting off again (2). At a first pass this run-and-tumble behavior is essentially diffusion, the bacteria covers a lot of ground without going anywhere on average unless there's some food nearby, in which case the run-and-tumble is slightly biased so that the bacteria drifts toward the food.

As I said, during runs the flagella are bundled together and all turning counter-clockwise. A tumble happens when the bacteria sends a signal to the flagella to start turning clockwise. This breaks apart the flagellar bundle and sends all the flagella in different directions. Up until the paper I mention above, it was still not clear exactly how many flagella need to start turning clockwise before the bundle breaks up and a tumble commences.

Now you can see why having more flagella might not necessarily be a good thing: you might be able to move faster during the run (though this is not obvious to me, given the sometimes unintuitive nature (3) of low-Reynolds number fluid dynamics) but at the same time the flagellar bundle will be more sensitive to break-up because with more flagella come more opportunities for one to start moving counter-clockwise and mess things up.

As the authors of the above paper (1) find, it turns out that there are correlations between the actions of different flagella in a bundle, so that a bacteria with many flagella effectively acts like a bacteria with fewer flagella. The result is that the overall pattern of run-and-tumble is not much affected by the number of flagella; you might argue that number of flagella has a minor effect on bacterial movement.

On the other hand, you might be interested not in the net movement of a bacteria, but just on its speed when it's moving forward, in the run phase. Unfortunately I wasn't able to find any good measurements of the run velocity of a bacterial cell as a function of number of flagella. But hopefully given this recent paper, maybe such measurements are on the horizon!

(1) Mears PJ et al. (2014) Escherichia coli swimming is robust against variations in flagellar number. eLife 3:e01916.

(2) Turner L et al. (2000). Real-time imaging of fluorescent flagellar filaments. J. Bacteriol. 182(10): 2793-2801

(3) Purcell EM. (1977). Life at low Reynolds number. Am. J. Phys. 45: 3-11.