What stops myosin during muscle relaxation? I understand that when the muscle is relaxed tropomyosin blocks myosin binding sites on actin filaments thereby preventing muscle contraction. What I am concerned with however, is whether myosin filaments - given that ATP is found in the muscle fiber- continue trying to bind to actin, and hydrolyzing ATP...
So how are mysoin filaments stopped from trying to bind to actin when the muscle is relaxed, to stop the hydrolysis of ATP?

Relaxation of vertebrate skeletal muscle is thought to occur in the absence of  ${Ca}^{2+}$ as a result of tropomyosin physically blocking the binding of myosin to actin. This steric blocking model of muscle relaxation predicts that myosin subfragment 1 (S-1) will not bind to actin under conditions where the acto-S-1 ATPase rate is inhibited.[1]
However, in a paper that questions this model, it is concluded that in the absence of ${Ca}^{2+}$, troponin-tropomyosin inhibits the ATPase activity by inhibiting a kinetic step in the cycle of ATP hydrolysis, perhaps $P_i$ release.[2]

That is the simple answer, but I think a little elaboration on the above paragraphs might help add some substance:[2]

Vertebrate skeletal muscle contraction is the result of a cyclic interaction of thick myosin filaments with the thin filaments, composed primarily of actin, troponin, and tropomyosin, causing these two sets of filaments to slide past each other. This cycling is driven by the hydrolysis of ATP by myosin in a reaction which is activated by actin.

When the sarcoplasmic reticulum lowers the free ${Ca}^{2+}$ concentration from $10^{-5}$ to $<10^{−7} m$, muscle contraction ceases and the associated actin-activated myosin ATPase activity is inhibited. The proteins troponin and tropomyosin are responsible for this effect of ${Ca}^{2+}$ on the interaction between myosin and actin.

At levels of ${Ca}^{2+}$ low enough to cause relaxation, tropomyosin is positioned away from the central groove of the F-actin filament where it appears that it might interfere with the binding of the myosin cross-bridge. This structural work formed the basis of the steric blocking hypothesis which suggests that relaxation occurs when tropomyosin, in the “relaxed” position, physically blocks the binding of the myosin cross-bridge to actin.

Three-dimensional reconstructions from electron micro-graphs have suggested that the myosin cross-bridge and the tropomyosin molecule may be in close contact with each other on the actin filament, a requirement for a steric blocking type model.

The steric blocking model predicts that in the absence of ${Ca}^{2+}$ the degree of association of S-11 should be much weaker with regulated actin than with unregulated actin. In fact, in the absence of ${Ca}^{2+}$, the binding of S-1·ADP to regulated actin is strongly inhibited in a cooperative manner.

At low levels of saturation of the actin filament with S-1·ADP, the binding of S-1·ADP to the regulated actin filament is about $10^3$ times  weaker than at high levels of saturation.

However, the fact that S-1·ADP binds weakly to regulated actin does not prove the steric blocking model since, in relaxed muscle, the cross-bridges normally exist with bound ATP (or ADP · Pi) and not bound ADP. Therefore, the steric blocking model predicts that troponin-tropomyosin should inhibit the binding of S-1·ADP·Pi as well as S-1·ADP to regulated actin in the absence of ${Ca}^{2+}$.

Using stopped flow turbidity measurements, we have previously measured the effect of ${Ca}^{2+}$ on the association of S-1 · ATP and S-1 · ADP · Pi with regulated actin. Surprisingly, in the absence of ${Ca}^{2+}$, the binding constant of S-1 · ATP or S-1 · ADP · Pi to regulated actin was only decreased to 56% of the value in the presence of ${Ca}^{2+}$ although the rate of ATP hydolysis under the same conditions was decreased to 6% of the rate with ${Ca}^{2+}$ present. These data suggest, in disagreement with the steric blocking model, that inhibition of the rate of ATP hydrolysis, in the absence of ${Ca}^{2+}$, is not the result of inhibition of the binding of S-1 to regulated actin.

1- https://pubmed.ncbi.nlm.nih.gov/27065174/
2- https://www.ncbi.nlm.nih.gov/pmc/articles/PMC1266292/#FN2


Trouble understanding overlapping genes. I read the wikipedia passage on overlapping genes.

https://en.wikipedia.org/wiki/Overlapping_gene

The Tandem overlaps they speak of make sense to me, as I understand that their are different reading frames(3 in total if I understand correctly) for genes to be read in the 3' to 5' direction and that this causes an out of phase overlap.

This out of phase tandem overlap for genes was shown by the picture below:

I am having trouble understanding the other scenarios spoken about in the passage, the in-phase tandem overlap and the Anti-parallel strand overlaps.
Maybe I just lack a good enough background in genetics.

I do not understand how you can have an in-phase tandem overlap. If the genes were in phase, how could they ever overlap and be different genes? Wouldn't they both respectively have a stop codon, and then one stop codon would "stop" the translation?

I picture it like this: pardon the crude drawing

The anti-parallel strand overlaps also confuse me as well, for I do not understand how genes can overlap when on opposite strands. If the genes are on opposite strands wont that mean they wont have any overlap in the first place? Am I missing something basic? Thanks.

A gene has got nothing to do with a stop codon. As you might know there are non-coding RNAs as well and the genomic regions responsible for transcription of these RNAs are also called genes.

Genes that have an overlap in their genomic locations (without considering the strand orientation) are called overlapping genes.

In case of same strand overlaps, the genes should have a different transcription start/stop site. Alternatively spliced RNAs are not considered overlapping transcripts. In certain cases, transcripts with an alternative transcription start/stop site are not considered separate genes but rather, isoforms, especially if the gene produces an mRNA and the alternate transcription start/stop sites just affect the untranslated regions (UTRs) and not the coding region (have a look at this article).

Your confusion is quite valid and gene nomenclatures need some redefinition/clarification.


Electrogenic ion transporter At the molecular level, electrical current across cell membranes ﬂows through three unique classes of integral membrane proteins (see Chapter 2): ion channels, electrogenic ion transporters, and electrogenic ion pumps.1

later it is said:

One class of electrogenic transporters includes the adenosine triphosphate (ATP)–dependent ion pumps.2

Is there any fundamental difference between "electrogenic" ion transporter and "electrogenic" ion pump, considering the fact that denosine triphosphate (ATP)–dependent ion pumps are also included in the former?

My attempt:
From this site ,Secondary active transporters are electrogenic, so if I consider those as electrogenic ion transporters , then also iam still confused why sodium potassium pump is in the category of ion transporter but not in ion pump.

Pump and transporter are synonymous: they refer to proteins that, unlike channels, use a form of active transport to move ions or molecules against their concentration gradient, although "pump" is most often used in the context of primary pumps, that is, those that directly use ATP, whereas "transporter" is more often used to describe the whole category or secondary pumps that use existing gradients.

  then also iam still confused why sodium potassium pump is in the category of ion transporter but not in ion pump

I have no idea where you are getting this statement from, clearly the sodium potassium pump is a pump.

Gadsby, D. C. (2009). Ion channels versus ion pumps: the principal difference, in principle. Nature reviews. Molecular cell biology, 10(5), 344.


What percentage of the genetic variance is explained by the 'n'- most important loci? Introduction
Standard models in population genetics look up at the evolution of few loci which impact fitness. The variance in fitness is determined by the genetic variance and the environmental variance (and the co-variance between environment and genetics). In this question I am interested only about genetic variance and about what percentage of the total (additive or not) genetic variance in fitness do 'n' loci explain.
Question
In general, in natural populations, what percentage of the total genetic variance is explained by the 'n'- most important loci?
Here, by "most important loci" I mean loci which variance explain much of the total genetic variance.
In other words, the subquestions are of the kind:

how much of the fitness variance does the most important locus explain?
How much of the fitness variance do the 3 most important loci explain?
How much of the fitness variance do the 100 most important loci explain?

Of course, the answer depends on the population under consideration. Factors that might influence the answers are for example

species
population size
environment stability
…

Beside this question, I also welcome some insights concerning how different factors are likely to influence the answer.

From the statistical point of view, this question is rather vague. One would need a mathematical definition for the term "genetic variance".

In one extreme, if the "genetic variance" merely means the categorial variations of nucleotides (i.e. ACTG) in the pooled genomes of interest, then the distribution of total "genetic variance" vs. loci variation is uniform and only depends on the size of the locus.

In another extreme (among many dimensions of extremes), if the "genetic variance" is only manifest by the organism's immediate "fitness" and only has two values: life and death (on birth), then all the "essential genes" are "the most important" loci. If you're interested in the n most important loci where n > the number of essential genes, then you would first look at the binary genetic interactions in the database such as BioGrid where two non-essential genes would "interact" and change the organism's fitness (in life and death).

Of course none of the two extremes is very interesting in population genetics or evolution, but a statistical question is best phrased by statistical terms. I would try to find the mathematical definition for "fitness variance", too.

---EDIT---

For a semi-empirical/informatics study, I think you could start with the simplest organism whose genome is well studied.

Choose an organism (e.g. yeast)
Assume uniform inheritability
Choose a specific measurable phenotype/environment (e.g. the ability to grow on a specific sugar x)
Scan each gene in the yeast genome and see its quantitative impact on growth (They're documented in various database)
Ignore genetic interaction
(Or scan each gene pair/triplet/.../n-cluster to see its impact on growth on x)
Try to model your empirical distribution. It's only valid for that specific phenotype/environment
Define your "TOTAL genetic variance in fitness" meaningfully and rigorously. "Additivity" would be a very drastic assumption.

My guess as a non-geneticist is that, as rg255 suggested, for each phenotype as a function of the environment, the distribution would follow a power law. They would not have the properties that would allow you to use central limit theorem to "add them up". But for a specific phenotype, your empirical cumulative distribution function (cdf) would answer your question.


Visual phototransduction, neurons depolarization when no light: why? While reading into visual phototransduction I was surprised to learn that photoreceptor cells are actually depolarized when there is NO excitation (no light, i.e. when you see nothing, black) and hyperpolarized when there is light (I expected the opposite).

This got me thinking...

Let the collection of neurons located within the eye be referred-to as IEN (inside-eye neurons) and, similarly, the collection of neurons directly linked to the eye, yet external to it be referred-to as OEN (outside-eye neurons).

Let the collection of neurons directly getting input from the retina be referred-to as RN (retinal neurons). I'm assuming RN < IEN (subset of -
 correct me if I'm meaningfully-wrong).

And finally, let the collection of neurons that directly link to OENs but are part of IEN be referred-to as XN (eXporter-neurons, in the sense that they export visual information from the eye to the "listeners" outside of it).

So the path of visual stimulus would then be RN -> [...] -> XN -> OEN.

Now, if my eyes were to be completely popped-out I'm assuming that relative to OENs that would have the same effect as XNs not passing along any excitatory/inhibitory information at all. Also, since what I'd perceive would be pure black, that seems to imply that XNs also don't get activated at all when we actually see black (e.g. close our eyes in a very dark room).

So, when we perceive black => RNs get depolarized => they release glutamate continuously => XNs don't excite/inhibit OENs at all (seems to also imply that RNs don't intersect XNs). It is this inversion that intrigues me - why is it that for XNs to transmit NO information RNs have to continuously be excited? By that sense, this makes me think of the eye as if it was a big NOT gate.

Short answer
The reason for the depolarized state of photoreceptors in the dark is unknown as far as I know. It probably has its roots deep in the evolution of the eye.

Background
Indeed, photoreceptors are continuously active in the dark. A steady current flows through open channels (mainly Na+ ions). This current is called the dark current and it partially depolarizes the photoreceptor cell. The depolarized photoreceptor releases neurotransmitter (glutamate) to secondary neurons. Upon illumination of the photoreceptors rhodopsin molecules are isomerized to the active form and the photoreceptor cascade (Fig. 1) is initiated. This ultimately closes the cation channels in the photoreceptor's membrane, stopping the dark current.  This leads to cell membrane hyperpolarization and inhibition of glutamate release to the second-order neuron (Fig. 1).

Admittedly, I do not know why phototransduction exactly is organized this way. There is a lot known about the evolutin of phototransduction (e.g., Lamb et al., 2016), but why the dark current exists is unknown as far as I know. I surmise it is probably an evolutionary relic. Evolution eventually has to work with pre-exisiting features to generate new structures, or alter them. Simple photosensitive spots evolved into the complex eyes observed in today's vertebrates. Evolution thereby used what was there and that happened to be an energetically unfavorable 'NOT' gate.

Fig. 1. Phototransduction in photoreceptors. source: Kolb, (2012)

Reference
- Kolb, Webvision. The Organization of the Retina and Visual System (2012)
- Lamb et al., Molec Biol Evol (2016); 33(8): 2064–87


And another insect from SE Brazil Oct 2017 Photographed in the low areas of the Brazilian Atlantic Rainforest - Oct 2017 and about 2cm in length.

I'm having a difficult time pin-pointing the species, however, I'm fairly certain this true bug belongs to the genus Zelurus.
Zelurus sp.

Zelurus sp.

Z. variegatus

Distribution

Interactive version can be found here.

source


what is the difference between osciloscope and electrocardiograph? I am completely confused about these two devices! as far as I know, oscilloscope and electrocardiograph, both, show the heat waves(P, QRS, T) so what is the difference between their job?

I would say that electrocardiograph is a specific case of oscilloscope. The main difference, perhaps is that "oscilloscope" is a device for visualization on a display ("scope" from greek "look out"), whereas "electrocardiograph" is primarily used to record ("-graph" from greek "to write").

Oscilloscope is design to visualize an electrical signal, that changes in time (generally speaking). Nature of that signal is not important, it can be radio signal, electrical signal from some device (Arduino), or sensor (thermometer). You can detect earthquakes if you have particular sensor for vibration.

It is just that these devices should output signals of particular frequency, usually general laboratory oscilloscopes are designed to work with signals up to few MHz. If you attach sensor that "feels" muscle contraction, you will have a tool for electromyography.

Electrocardiograph, by definition, is a device that visualizes electrical signals acquired via electrocardiography. It is just that usually when people and doctors perform electrocardiography, they use specialized devices, not general-purpose oscilloscope. For example, cardiac signals are usually well below 0.5-1KHz, you can save a lot of money and complexity by not having to worry about higher frequencies in your device.


Why can't complex multicellular organisms reproduce by regeneration? My textbook says that regeneration isn't  possible in complex multicellular organisms.

But why? Why can't specialized cells be present in the bodies of complex multi cellular organisms which can proliferate and ultimately make a whole new organism, when (for example) cut in two?

I was doing some research on this myself and from what I had found, the body has hox genes that direct the formation of body parts. So now the obvious question is that if a body part contains stem cells and hox genes, shouldn't it be able to regenerate a new body? Well, perhaps. But the problem is that as organisms become more complex, the process of the creation of their bodies also becomes more complex. For example, in the fruitfly, if the labial gene function is not expressed, the mouth and head that develop outside the body do not get involuted and that prevents the salivary gland from forming.

So the answer about why multicellular organisms can't regenerate so easily simply seems to be because the process of creating a multicellular body has for some reason become an extremely complex, intertwined process where various body parts creation is dependent on other body parts creation. So in the evolutionary process it's entirely possible that the stem cells noted that there's no point trying to regenerate a body part, and therefore they just don't.

Here's an anomaly. There's a jellyfish that can reverse aging and go back to infant stage. I think it's able to do so because it has a simple body structure.

Another possible reason could possibly be the presence of unstructured proteins which are more abundant in multicellular organisms.


Evolutionarily, why do male rats and horses lack nipples? Developmentally male rats don't have nipples because (reddit)

  Testosterone release in the fetal male rat happens before the stage of mammogenesis where the teat is formed whereas other species halt that same process after the teat is formed.

According to the same page, male rats and horses (and mice) are unique among non-monotreme mammals in not having nipples.

Male humans have nipples because there's no evolutionary pressure for them not to. (Why do men have nipples?) Then, why was there selection for males to lack nipples in the case of mice, rats, and horses?

(Apparently, male horses lack nipples because "female horses have their nipples between their hind legs, and in males, that real estate is used for other purposes". (Straight Dope))

Evolutionarily, why do male rats and horses lack nipples?

Male mice lack nipples too. Mice are frequently used for embryonic research as they are small and reproduce quickly. It is thought that male mice do develop nipples, but that they regress during development (Wysolmerski, 1998).

In general, it is thought that mammalian organisms develop as females by default when there is no male (Y) chromosome present (Hughes, 2004). Notably, X_ individuals and XY individuals with certain Y-chromosome deletions also develop as females. Basically, it is testosterone that suppresses female features. Hence, also boys with testosterone production, but without the appropriate receptors also develop as females.

Hence, men have nipples because females do (Simons, 2003).

The development of nipples is a pretty complex process (Wysolmerski, 1998). Disruptions in the process cause regression. Indeed, in mice, the tissues typical for nipples do develop in males during embryogenesis, but degenerate within a few days, leaving no trace of nipples at birth.

Hence, given the mouse example, I think that every mammalian male species develops by default nipples somewhere in development like mice, but that development is halted at various stages across species. The later in development it halts, the more it resembles female nipples. When it is stopped very early after gestation, they regress completely, like in mice.

Why is there a difference between species as to when the nipple development halts? Because the nipple developmental sequence is complex (Wysolmerski, 1998), and because evolution typically prefers losing things over gaining them, it is likely that by happenstance one of many crucial genes was lost in males, but selective pressure was simply too small to be of any significant advantage over males with slightly better-developed nipples. But from here on, it is guessing.

References
- Hughes, N Engl J Med (2004); 351(8): 248-50
- Lawrence, Nature News, August 1999
- Simons, Sci Am, September 2003
- Wysolmerski et al., Development (1998); 125: 1285-94


Can lymph be in peripheral blood? I read an argument that 1-3% of lymph is in peripheral blood.
However, I am not sure if this lymph is about lymphocytes in peripheral blood; not lymph itself.
Lymph gets exchanged between capillaries through Starling law powered by the heart and Frank-Starling mechanism.
In this way, it makes sense that some lymph is leaking out from capillaries into the bolod.

Can lymph be in the peripheral blood?

Depends if your talking about interstitial fluid or lymph:

As per starling forces: the oncotic pressure of the late capillary plus hydrostatic pressure of interstitial space, would be higher than oncotic pressure of interstitium [though I dunno why since intertitial GAGs are very water hungry and they are responsible for the negative (hydrostatic?) pressures found in interstitium... Read lots of articles, still makes no sense to me] and hydrostatic pressure of the capillary... Thus there would be a net influx of interstitial fluid in the very ending part of the capillary, toward the venous side...

I say interstitial fluid, because it's not lymph yet since the interstitial fluid hasn't yet entered the adjacent lymph capillaries... if the interstitial fluid entered the lymph capillary then it would be stuck there, the mechanism of this one way flow is the junctions between the endothelial cells in these blind ended lymph capillary... These junctions which only opens up gaps between these endothelial cells when the interstitial pressure outside the lymph capillaries is higher than pressure inside the lymph capillaries...

The lymph then goes forward through lymphatic vessels/nodes because external pressure (pulsing of adjacent artery, contraction of adjacent skeletal muscle, changes in thoracic pressuring during breathing, etc) and there is little back flow because of valves (just like venous valves which prevent back flow in the low pressure venous system)... Ultimately the lymph returns by either thoracic duct (in the angle between left subclavian and left jugular vein) or right lymphatic duct  (in the angle between right subclavian and right jugular vein)...

There are some sweet review articles from University of Bergen, Norway that go much more in depth on the physiology of interstitial fluid...
Here is one of their more recent review articles:
http://m.physrev.physiology.org/content/92/3/1005 it's free

IN SUMMARY: the starling forces bring interstitial fluid into the blood (the "peripheral part" right after the capillary)... The thoracic and right lymphatic duct drain lymph into the blood (I'm assuming this does not apply by your definition of peripheral, since this venous blood is about to be drained into the right atrium and is destined for the lungs)...

P.S. aqueous humor (in the eye), endolymph (cochlea aka inner ear) and CSF are also absorbed back into the blood... Also the brain has an interstitial fluid that is different from the CSF found in the subarachnoid spce of the meninges, i actually don't know how that drains... But anyway, the head is so special, I'm sure you weren't wondering about these specifics...


Protein function and how it does not get affected with amino acid How is it possible that a protein's function is not changed when an amino acid within it is replaced?

The side chains of some amino acids are not reactive and therefore not involved in any covalent chemistry in enzyme active centers.

Isoleucine, an essential amino acid, is one of the three amino acids having branched hydrocarbon side chains. It is usually interchangeable with leucine and occasionally with valine in proteins.

Typical amino acids and their alternatives usually have similar physicochemical properties.

Amino acid replacement is a change from one amino acid to a different amino acid in a protein due to point mutation in the corresponding DNA sequence. It is caused by nonsynonymous missense mutation which changes the codon sequence to code for another amino acid instead of the original.

The effects of amino acid replacement may vary depending on how similar or dissimilar the replaced amino acids are, as well as on their position in the sequence or the structure.

References:

http://www.biology.arizona.edu/biochemistry/problem_sets/aa/Isoleusine.html
https://en.wikipedia.org/wiki/Amino_acid_replacement


JMol "calculate HBONDS": which atom is the donor/acceptor? JMol can be used to identify Hydrogen bonds in proteins by the script "calculate HBONDS".  By outputting the state of the network, we can obtain a lisdt of H-bonds.

Here's one line from an example output (PDB ID number 1a1x):

74    540    4096    0.001    -2.507   hbond;

I believe the correct interpretation of the first two numbers is that atom 74 is H-bonded to atom 540.  (I'm not really sure what the other numbers mean.)

  Question: How can we determine which atom is the donor and which is the acceptor?

In the above example, atom 74 is a Nitrogen atom and atom 540 is a Carbon atom, so presumably the Nitrogen atom is the donor.  But, how can we tell in general?

I dont know about JMol, but this can be easily done with UCSF Chimera.

I loaded the 1a1x structure.
Then selected: Tools > Structure Analysis > FindHBond

I kept the default values and selected Write information to reply log

Then clicked OK

On the structure the H Bonds will be identified with a coloured line.

To view which is the donor and acceptor go to: Favorites > Reply log

An output of the reply log is below:

H-bonds (donor, acceptor, hydrogen, D..A dist, D-H..A dist):
ASP 12.A N    GLU 25.A OE2  no hydrogen  3.012  N/A
HIS 13.A N    GLU 25.A OE1  no hydrogen  3.338  N/A
HIS 13.A ND1  ASP 12.A OD2  no hydrogen  2.835  N/A
LEU 14.A N    TRP 70.A O    no hydrogen  2.934  N/A
TRP 15.A N    ARG 23.A O    no hydrogen  2.866  N/A
...
120 hydrogen bonds found

As you can see the donor and acceptor has been identified as defined in the header line.

If you wish, the information can also be written to a file by selecteing Write to text file in the FindHBond window.


Difference between morphogens and evocators in development? I am specifically interested in this within the context of 1950's developmental biology (I presume the meaning and scope of usage of the terms has changed since then), as I am reading Alan Turing's paper on 'The chemical basis of Morphogenesis'. Turing first mentions evocators in the context of 'evocators of Waddington'. I have read some articles by Waddinton to try to see what he meant by 'evocator', but have not gained clarification from the articles available online. Turing cites Waddington's "Organisers and Genes", which is not available on the web and I do not have access to.

A more recent article here states

  The first element is the inducer (or evocator) defined as the signal emanating from a piece of tissue, such that when the tissue is transplanted to an ectopic position it elicits changes in the behavior of neighboring cells.13 Today, we understand that most of these inducers are secreted proteins and include morphogens belonging to the Wnt, Hedgehog (Hh), TGFβ, EGF, and FGF families.

Which seems to imply that most evocators are morphogens, although from the definition of evocator given I would say that this is exactly what a morphogen is!

Any insights would be much appreciated!

"Since the evocator Waddington and Needham discovered was ether soluble, it was considered to be a sterol, a class of molecules that includes many hormones related to later development" (https://embryo.asu.edu/pages/chemical-induction).

It seems like the termed is outdated, hopefully, you find this information useful.


Problem on number of amino acids that an alien can use A new life form discovered on a distant planet has a genetic code consisting of five unique nucleotides and only one stop codon. If each codon has four bases,what is the maximum number of unique amino acids this life form can use ?

I tried this and got 624 (Using permutation method 5×5×5×5 = 625). But the answer given was 3124. Please help.

I also got 624. Maybe the answer is incorrect.
My processing is:

Each codon has four bases, the species contain five unique nucleotides, so for the first position in the codon, there are 5 possible results, also, for the second position in the codon, there are 5 possible results, it is same in third and forth position,

so the result should be  5*5*5*5-1=624


Is RNA polymerase affected by proteins bound to the coding sequence of a gene? I am designing a synthetic gene construct to express genes in E. coli driven by either Ptet or PLacO. The construct would look like:

-Ptet-(Gene1)-PLacO-(Gene2)-

I want to express each gene using either aTc or IPTG, but I want to make sure that the transcription of each gene can be controlled independently.

Let say I only add aTc and transcription is initiated at the Ptet location, would RNA polymerase bump into LacI bound to the PLacO promoter and stop transcribing? Or would it knock-off LacI transcribe (Gene 2)? I am not sure whether or not need to include terminator sequence to make sure the expression of each gene is decoupled.

Nice question!
First of all make sure that you have multiple transcription stop / terminator sequences at the end of your first gene . This is quite standard procedure. Also in your case the phenomenon called transcription interference (TI) comes handy. Here's a review on TI. In short:

  the term TI usually refers to the direct negative impact of one transcriptional activity on a second transcriptional activity in cis.

Although TI is mostly defined as two transcription events affecting each other, at the low level the basis of it is that two proteins bound to the DNA affect each other's activity. Now since LacI's function is to inhibit transcription, thus it is strongly bound to the DNA and I'd say even if a polymerase somehow slips through the polyA signals at the end of the first gene it is likely to be stopped by the LacI already bound at the second promoter.

Alternatively you could put a weak promoter opposing the first gene and drive its expression, so it would collide with the oncoming polymerase from the first gene thus stopping it before reaching the second gene. Also here is another article on TI, that proposes a gene regulation network based on TI.


What is meant by translational repression by an upstream open reading frame? The NCBI report on the human huntingtin gene contains the following statement:

  “This gene contains an upstream open reading frame in the 5' UTR that
  inhibits expression of the huntingtin gene product through
  translational repression.”

But no specific reference it given. Can anyone explain what this means.

Summary

The statement refers to the small but significant translation of short open reading frames (ORFs) of huntingtin mRNA, 5′ to the main open reading frame that encodes the large (3144 amino acid) functional product of the mRNA (and gene). It has been reported that the translation of these short ORFs decreases (i.e. inhibits) the translation from the main ORF. It is implied that this is a mechanism for regulating the expression of the gene, and, as it occurs on the mRNA in the cytoplasm, it is an example of ‘translational regulation’ — specifically inhibition. The precise mechanism of this regulation — although not unique to this mRNA — is apparently unknown.

Translation of ORFs in eukaryotic mRNA

Most eukaryotic mRNAs are functionally monocistronic — i.e. they encode only a single protein — unlike bacterial mRNAs. Initiation of protein synthesis involves selection of an initiation codon (usually AUG) by a complex of initiator-tRNA and the small ribosome subunit, but, in contrast to prokaryotes, the mechanism in eukaryotes involves scanning along the mRNA from the 5′ end. In the majority of cases the initiation codon of the protein product is the first AUG, but the significant number of exceptions suggest an effect of sequence context or mRNA secondary structure.

In recent years evidence has accumulated that in cases where the first AUG is not that of the main protein product, a small amount of translation from such 5′ upstream ORFs does often occur.

Diagrammatical representation of mRNA with upstream ORF (blue with initiation codon represented as an open triangle) and main ORF (black with initiation codon as closed triangle.)

Regulatory role of upstream ORFs

Johnstone et al. (2016) have reviewed a large body of evidence that indicates that upstream ORFs are translational repressors in eukaryotes. This has generally been established by measuring the extent of translation of the main ORF after manipulation or removal of the upstream ORFs. Much less is known about the actually mechanism of such regulation, which may vary from case to case, and certainly the paper on huntingtin only offered speculative possibilities. These include:

Binding of small peptide products of the upstream ORFs to the AUG of the main ORF
Pausing of ribosomal initiation complex at upstream ORFs impeding the progress of other to the AUG of the main ORF
Change of mRNA secondary structure by the binding of the ribosomal initiation complex at upstream ORFs increasing the difficulty of the initiation step in which the 5′ end is ’melted‘.

It also appears that the extent of the translation of upstream ORFs — and hence the extent of the repression — can itself be regulated by mechanisms such as the phosphorylation of initiation factor eIF2.


What are some places where biofilms could develop? I'm trying to think of places where a biofilm could develop other than on medical equipment or food processing equipment such as stainless steel mechanized blades or knives. I'm thinking more along the lines of biofilms on typical house hold things such as the bottoms of shower curtains

They can certainly develop in the toilet. Just try not to disinfect for a week...

Another place is the surface of living organisms, e.g. dental plaque or the leaves of a plant, etc...


Phagocytosis and bacteria Invading Bacteria are phagocytosed by phagocytes. Phagocytes are specialised white cells. But, if a non-phagocyte cells are infected with bacteria how those cells handle this situation? Do they succumb to the bacteria?

For cells to become infected with bacteria, some bacteria will still be present in the blood. These are still engulfed by phagocytes in the same fashion. However once engulfed the antigens on the surface of the Bacteria are placed on the surface of the phagocyte, this is known as an antigen presenting cell (APC). Specialized lymphocytes known as Th Cells bind to these antigens and then divide. These then stimulate Cytotoxic T cells which secrete a chemical known as perforin. The perforin makes holes on the membrane of 'non-phagocyte' infected cells allowing water to enter. This essentially causes the cell and the bacteria to explode, known as cell lysis.

So basically the infected cell is simply destroyed unless hydrolytic enzymes break down the Bacteria.


How to determine if the amount of manganese chloride will change neural resting potential If I am treating an organism in MnCl2 dissolved in water, how do I determine if the amount of Cl (in MnCl2) will change the neural resting potential Oor if it will influence motoneurons?

For theoretical calculations with other ions, you need to know the concentration and permeability of each of the major ions inside and outside the cell, both with and without the added substance (in your case MnCl2), and then use the Goldman equation/GHK equation to calculate resting potential before and after. Approximations would probably be sufficient.

However, chloride can be a bit complicated because the intracellular chloride concentration can fluctuate with extracellular chloride moreso than the other major ions. Therefore, depending on the sensitivity of your experiment, you may need to perform whole-cell patch clamp recordings to measure the resting potential directly.

If the change in chloride is a small percentage of the total, you can probably consider the influence to be negligible.


How much do crocodiles eat? How much do crocodiles eat? The closest I found to an answer was on German wikipedia, where it says that adult nile crocodiles eat "only about 50 full meals a year." While interesting, without knowing what a full meal is this is meaningless. So how much does an adult crocodile actually eat?

I think relevant other questions in this context are:

Do all adult crocodiles (or even all crocodilia) eat roughly the same (corrected for body mass)?
Some crocodiles apparantly eat fruit and seeds on a regular basis*, is this is a relevant part of their caloric intake or just a "vitamin supplement"?

*paywalled article, I only read the abstract

It's only a partial answer, but a Nile crocodile can eat up to half its body weight at a sitting, according to National Geographic. Since few freshwater fish weigh half as much as a crocodile, we can probably assume these croc feasts occur when they're able to ambush ungulates (e.g. zebras and antelope).

So one guess is that a crocodile feasts on ungulates about fifty times a year, though its normal diet is fish. Of course, I could be wrong - perhaps five fish caught in the space of an hour are considered a full meal as well.

I never knew crocodilians ate fruit. (The link you provided doesn't mention fruit, by the way.) But this article says they can eat "large quantities" of fruit, suggesting they're looking for more than supplemental vitamins. But who knows? It sounds like scientists don't yet know much about fruit-eating crocodilians.


Can plants grow without photosynthesis? Animals can survive without eating if nutrition is injected directly into their blood.

Can the equivalent be done with plants? By injecting carbohydrate rich nutrient solution, either without photosynthesis or in addition to photosynthesis.

To answer if the equivalent could be done with plants, we have to understand the histological and anatomical differences between plants and animals.

Just as spraying a nutritive solution over an animal (like you and me) doesn't work, spraying a nutritive solution over the leaves doesn't work as well: the nutritive molecules (let's use glucose in our hypothetical example) don't penetrate the leaf's epidermis. The same happens if you put the nutritive solution in the soil: there are several membrane layers between the external medium and the phloem, and glucose will never reach its destination.

Now, let's try to compare with a IV injection of nutrients, which we can do in an animal. The plant's structure that (kind of) can be compared to blood vessels is the phloem (actually, xylem and phloem, but in our case, dealing with nutrients, we're gonna talk only about the phloem). To inject the nutrients we have to reach the phloem, because the periderm in plants with secondary growth ("cork" in the image) is very impermeable:

Now, even if we reach the phloem, using methods like this...

...we still have a problem: in an animal like you and me, a substance injected in a given vein will quickly spread to all the systemic blood vessels. But the same doesn't happen with the phloem. The injected substance can go up or down (depending on the pressure flow), but it doesn't spread to all of the plant's tissues, unless we perforate the phloem in a lot of places.

Thus, so far, this method is used to deliver chemicals (like pesticides) into the tree. Notice that, just like a person in a hospital, the bag with the solution has to he higher than the perforation:

For a chemical treatment like this, the amounts of the substance are small. But for "feeding" a plant using this technique, the amounts of necessary nutrients are way higher...

Finally, answering your question: is it hypothetically possible to feed a plant in the dark with a nutritive solution? yes. Is it practical, or doable? very hardly.

Source: Utah State University


Why does frost increase soil fertility? I was reading this paper: https://link.springer.com/article/10.1023/A:1011398431524, which makes claims based on the the idea that frost can increase soil fertility. How/Why will frost increase the fertility of the soil? I was always under the assumption that frost was bad for soil and plants in general.

I found this article on why frost increase soil fertility - specifically nutrient density in the soil.  101 years old!
The Occurrence of Bacteria in Frozen Soil
E. C. Harder  Botanical Gazette  Vol. 61, No. 6 (Jun., 1916), pp. 507-517
https://www.jstor.org/stable/2469064?seq=3#fndtn-page_scan_tab_contents
Relevant text pasted below.  In short: it is bacterial action and nitrogen fixation which contributes to soil fertility.  Freezing increases bacterial population and these metabolic activites.  The theory is that freezing kills bacterial predators on higher trophic levels (protozoans) but the bacteria live and thrive in the absence of predators.

An analogous hypothesis proposed by RUSSEL3 for increases in the number of bacteria after partial sterilization by heat, frost, or other means is that by such partial sterilization the protozoa are killed, thus permitting the unhindered development of bacteria which under normal conditions is held in check by protozoa.
BROWN and SMITH (loc. cit.) in their investigations dealt mainly with the physiological activities of bacteria under conditions of low temperature and frost, although they also made some determinations of the number of bacteria in frozen soil. Their principal conclusions regarding the ammonifying, nitrifying, denitrifying, and nitrogen fixing powers of frozen soils are as follows: (1) that "frozen soils possess a much greater ammonifying power than unfrozen soils"; (2) that "during the fall season, the ammonifying power of the soil increases until the temperature of the soil almost reaches zero, when a decrease occurs, and this is followed by a gradual increase and the ammonifying power of the soil reaches a maximum at the end of the frozen period"; (3) that "the nitrifying power of frozen soils is weak and shows no tendency to increase with extension of the frozen period"; (4) that "frozen soils possess a decided denitrifying power which seems to diminish with the continuance of the frozen period"; (5) that "during the fall season, the denitrifying power of the soil increases until the soil freezes, after which a decrease occurs"; (6) that "frozen soils possess a nitrogen fixing power which increases with the continuance of the frozen period, being independent of moderate changes in the moisture conditions, but restricted by large decreases in moisture"; and (7) that "in the fall, the nitrogen fixing power of the soil increases until the soil becomes frozen, which in almost ceases, after which a smaller nitrogen fixing power is established."


Blood pressure in various blood vessels and heart What are the typical systolic and diastolic blood pressure in arteries,capillaries, veins and the heart?

I know it must be highest in the arteries. arteries>capillaries>veins

Arteries: 120-80mmHg (systolic - diastolic). I also found mean values for capillaries (30mmHg) and veins (10mmHg) but not able to find systolic and diastolic pressure specifically.

Blood vessels have "compliance": that is, they are stretchy. Mathematically, this is very similar to capacitance in electricity.

The further from the heart you go into the vascular tree, the more the oscillations are dampened and the less it makes sense to think of systolic and diastolic pressures. In the capillaries and veins, blood flow and pressure is more or less constant rather than oscillatory (except in the most proximal veins where the motion of the heart again adds some pulsatile influences).

Therefore, for those vessels, you can just stick to mean values for pressure, the oscillations around those values are small.

See another Q&A: Why do people perceive blood pressure as the force that moves the blood forwards (see details)? that has a diagram showing the dampening of oscillations in pressure.


Alkaline Phosphatase and Ligase Protocol for Cloning In the image the circular molecule is a restricted vector and the linear red molecule is a DNA insert.
I found this protocol in my lessons notes, but I don't understand how it is possible that ligase catalyzes a reaction between 5'-OH and 5'-P. This reaction should form as product a 5'-5' bridge between the 2 nucleotides of the insert and the vector.
Usually Ligase repairs nicks between 5'-P and 3'-OH thus creating 5'-3' bonds between nucleotides that are normal in nucleic acids, so my question is how is it possible for the DNA to be functional even though it has a 5'-5' bond?

DNA ligase catalyzes the formation of a bond between 3’-OH and 5’-PO4 of its nucleic acid substrate. (See Lehnman, 1974) The ligation bond will be formed when the proper ends are close together, allowing the ligase to act. Lets redraw your diagram and add the 5’ and 3’ designations to it. I think this may help clear out some points.

Regarding the alkaline phosphatase treatment with ligation, I think Dale and Greenaway said it so much better than I can. Quoting (Dale and Greenaway, 1985):

[The strategy] is to treat the cleaved vector DNA with alkaline phosphatase that removes the 5’-phosphate group from the DNA. Since the presence of a terminal 5’-phosphate is necessary for the action of DNA ligase, a dephosphorylated DNA strand will not be able to participate in a ligation reaction. Since this will apply to both ends of the linearized vector, the ends cannot be joined together again, i e , parental vector regeneration is impossible. This results in a very low frequency of transformation by the vector (since linear DNA transforms very poorly). However, the vector 3’-hydroxyl ends are available for joining (by DNA ligase) to the 5’-phosphate termini of any DNA fragments to be cloned. This will result in a circular recombinant molecules with a single-stranded nick at each end of the vector portion. These molecules usually transform as efficiently as the fully ligated product, and the nicks are repaired in vivo after transformation. Thus the phosphatase treatment of the vector molecules increases the proportion of hybrid DNA molecules in the population of transformed cells, and also prevents the formation of oligomeric forms of the vector.

References
Dale, J.W. and Greenaway, P.J. (1985) “The use of alkaline phosphatase to prevent vector regeneration” Methods Mol. Biol. 2:231-236. Link to paper
Lehnman, I.R. (1974) “DNA Ligase: Structure, Mechanism, and Function” Science 186:790-797. Link to paper


Why do the brains of cocaine-users shrink faster than the brains of non-cocaine users? http://news.sciencemag.org/sciencenow/2012/04/cocaine-may-age-the-brain.html?rss=1

  Cocaine-dependent individuals showed a significantly greater-than-normal age-related decline in gray matter in prefrontal and temporal regions compared with healthy controls. By contrast, parts of the striatum appeared resistant to age-related volume decline in the cocaine-using group. Enlarged striatal volume has frequently been reported in stimulant-dependent individuals,5, 6 possibly reflecting a marker of reduced dopamine neurotransmission in this dopamine-rich brain region where drugs like cocaine work. Decline in striatal dopamine receptor density has been associated with normal age-related cognitive decline.7 The relative absence of age-related changes in the striatum of cocaine-dependent people may thus reflect another feature of an abnormal brain ageing process.

Could it be related to possible neurotoxicity, or to the long-term effects of dopamine receptor downregulation? Could the effects also be generalizable to those who take dopamine-based ADD medications?

A link to the Letter to the Editor appeared in Molecular Psychiatry to which the ScienceNow article refers:
Cocaine dependence: a fast-track for brain ageing?

The mechanism of action of cocaine is dependent on pre-existing dopamine production and secretion. Normally, secreted dopamine is cleared from the synapse via the dopamine transporter (DAT) located on presynaptic dopaminergic terminals. Cocaine inhibits
this reuptake of dopamine, increasing it's duration of action on post-synaptic dopamine receptors. Thus, cocaine's effect depends on neuronal dopamine production. Without intrinsic dopamine, cocaine would have no effect. This is actual not exactly true though because cocaine also inhibits serotonin and norepinepherine reuptake transporters as well but the same argument applies to those neurotransmitters as well.

The overall effect of cocaine on the nervous system is extremely complex as it prolongs the action of dopamine, serotonin and norepinephrine wherever those reuptake transporters are present. There is a particularly high density of DAT in the basal ganglia and nucleus accumbens. What these areas of the brain normally do is under intense study and debate. The accumbens may play a role in goal directed behavior and incentive salience, which is the attribution of value to various actions or objects in the environment. This is what naturally guides us to perform one action over another. The simplistic view is that cocaine "highjacks" this system such that cocaine itself acquires greater incentive salience than natural reinforcers such as food, sex, money etc.

Getting back to the original question regarding why cocaine causes brain atrophy. This appears to be a finding quoted from a paper by the original inquirer. It is not obvious to me what the mechanism of this atrophy would be. I suggest you look at the discussion section of the paper from which that abstract was quoted to see what the authors suggest. However, to my knowledge, such atrophy has not been widely discussed in the cocaine literature but I could be wrong.

For further information see Sulzer. 2011. How Addictive Drugs Disrupt Presynaptic Dopamine Neurotransmission.


How does the MET gene work and what happens when the promoter region gets mutated? I am doing research on inherited risk of Autism Spectrum Disorders(ASD) due to common Copy Number Variants(CNVs) One of the mutations is the 'CC' variant of Rs1858830 in the promoter region of the MET gene. (http://www.snpedia.com/index.php/Rs1858830)

After digging a bit deeper I found that hepatocyte growth factor is a pretty ubiquitous compound in the body and it's necessary for growth/repair of many tissues, namely the liver and digestive tract. I am trying to separate the direct causal effect of the mutation itself from the effects it has on other tissue which may themselves have an effect on ASD risk.

My question comes from a study which measured reduced serum HGF levels in autistic children with this mutation. (http://www.ncbi.nlm.nih.gov/pmc/articles/PMC3694825/)

I am confused about the implied causal effect of the 'CC' variant on HGF levels even though it is merely the ligand for the MET gene? Please help me understand where in the gene signalling process does HGF come in, what does it do, what proteins are transcribed by the MET gene, and how this mutation effects the transcription of these proteins.

HGF is the ligand for Met, which is a membrane receptor involved in multiple transduction pathways. Gene ID 4233 we're looking at a single gene product for MET.

If we're just reading through the wikipedia page for c-Met, when HGF binds to Met, it activates the tyrosine kinase activity of Met. Met recruits Gab1, and this complex mediates interactions with Ras, P13K, ß-catenin.. not just pathways associated w/ Autism, but also associated with invasive growth cancers. An interesting note about the C allele, it's associated with decreased Met transcription. Quoting a recent article,

  The present morphological and functional studies have revealed a highly novel mechanism conferred by MET receptor signaling, which we show serves a pleiotrophic role in controlling both neuronal and spine morphology, and the time course of glutamatergic synapse maturation on CA1 hippocampal neurons. Such findings are consistent with the possibility that mistimed maturation of glutamatergic synapses underlie aberrant function of neural circuits that are enriched in MET during development.

Source

So the critical results from the same study showed that when Met got knocked down or deleted, we saw maturation of neural tissue. So now, if i'm thinking about Met and HGF, your article notes that HGF is decreased in autism brains. If we (1) need HGF for Met function, and (2) ASD may be attributed to decreased Met function, and (3) the C allele decreases Met levels, we're basically looking at a system that's modulating for ASD.


MS/MS for complex samples I saw this figure in this atricle:

It's a really nice representation of what happens when using MS/MS. However I was wondering how this works when there are more proteins present in the gel spot. How would you be able to select the parent ion(because there are multiple of them)?

Liquid chromatography (LC) allows the peptides from a complex sample to be fed into the mass spectrometer over time. The bottom half of the image shows an MS spectrum from a single precursor scan and the machine is programmed to select the three most intense peaks (parent ions or precursor ions) from this scan for MS/MS (MS2).

This process is repeated throughout the LC-MS/MS run and programs usually include exclusion criteria where a parent ion is ignored for 30 - 90 seconds once it has been scanned 1 - 2 times. Usually the machine will also ignore singly charged parent ions as many of these ions will be inorganic contaminants and you want the machine to be busy with peptides.

Efficient use of machine time is critical to getting good coverage of the proteins in the sample. Proteins that are small create fewer peptides that are amenable to sequencing so these can be missed. Also those proteins that are less abundant may produce peptides that stay below the selection criteria. Good chromatography and efficient use of machine time are critical to getting good results and even in good conditions most peptides will not be subjected to MS2 sequencing. This estimates that of 100 000 peptides with only 16% were 'sequenced'. This explains why proteomics labs are speed freaks, and has driven the development of faster machines. The ~ 1 hour yeast proteome paper demonstrates the 'work' that the newer machines can get through;

  Over a 1.3 h chromatographic method .... collected an average of 13,447 MS1 and 80,460 MS2 scans (per run) to produce 43,400 peptide spectral matches and 34,255 peptides with unique amino acid sequences (1% false discovery rate (FDR)).


How long does it take for memory to disperse after death? 1: There seem to be cases where coma patients with a non-active brain (i.e. flat EEG) have regained full consciousness. => Apparently memory and knowledge are stored independent of brain activity.

2: There seem to be animals (e.g. hamsters) that can be frozen to complete organic inactivity and will regain full functionality after being thawed. => Apparently the stored memory does not depend on blood flow and other support.

3: From this I assume that quickly cooling a human brain to a temperature low enough to avoid decomposition would preserve the state of that brain that corresponds to that human's memories, knowledge, cogntivie abilities, and maybe consciousness at the time of cooling. => With the proper technology that "content" is theoretically retrievable.

Q: How long would that state remain after death in a brain left at room temperature?

Or in other words: How long does it take for decomposition to destroy memory?

Addendum:

The fact that the stored memory may not be accessible with current means is not relevant to my question. We cannot access the information stored in Linear A, but this unretrievability does not delete the information.

There are multiple levels of memory, some of which would die immediately, some of which would take some time. So the answer is: it depends; some immediately, some only very slowly.

At the highest level, the current neuronal firing state of the brain encodes memory on a very short scale - working memory. The memory held on this level does not have a clear anatomical counterpart (but for the potential encoded in the synapses). It equals very short-term memory/STM sequences, such as the words you read just before you read the words you're reading right now. This memory is lost immediately when you lose consciousness, at least to some degree; as this memory is hard to even strictly distinguish from consciousness and attention (though see Jonides et al. 2008).

Other forms of short-term memory/STM are stored in a slightly different form: short-term potentation, the adaption of neuronal responses following brief and intense stimulation. Spike Frequency Adaption/SFA is at an intermediate stage between this and the previous level. Short-term potentation and SFA decay within minutes or even seconds if they are not transferred into some more durable form of memory.

Long-term memory (/LTM) stores have specific anatomical correlates; they are stored in, amongst others, the synaptic weights (i.e. the amount of influence the firing of one neuron has on another). Some forms of LTM are best located in cortical synapses, others in the hippocampus. An even more fundamental, long-term, durable storage form is the wiring itself; not just the weights, but the existence of a synapse between two points, or not. For example, an important part of early learning is synaptic pruning, where synapses which do not play a meaningful role die off, whereas those which connect functionally related brain areas remain.
This pruning instantiates one form of learning, and the non-existence of a synapse is a form of memory. Synapses are comparatively stable. Even if the corresponding neurons die, in principle, the synapses still exist - and more importantly, the nonexistence of a synapse is even more durable. This form of memory can be observed in slice preparations of animals long dead.

For a simple example, consider any experiment on the neuronal responses in slice preparations, which can be considered a form of (decontextualised) memory access.

On the most extreme end, epigenetic adaptions and large-scale brain anatomy (which shows developmental traces) can be considered a form of memory that will remain intact until the whole structure rots away.

However, the more direct answer to the question the OP is asking is that once a large amount of neurons have died (brain death), there is currently no power on earth that can access a non-trivial amount of memory. As long as this is averted, nontrivial amounts of memory can be recovered.
For the parts that are lost immediately, see the source by Jonides et al. For more durable memory, you could look at for example Purves et al, Neuroscience.


Can a chicken's lard become liquid depending on its diet? On Cooking.se, somebody asked a question about the fat in chicken stock sometimes being liquid and sometimes solid once the stock has cooled. The only answer until now suggests that the difference is due to the percentage of saturated fat in the chicken lard, which will vary depending on the chicken's food.

I am not entirely happy with this explanation. I am aware that the proportion of saturated to unsaturated fat in an organism's cell walls will vary with diet, and could imagine that the fat stored within its adipose tissue will be different too. The person who provided the answer also posted references which confirm that part. But I doubt that the effect will be strong enough that, if an animal on low-saturated-fat diet is processed to lard, the resulting lard will be a viscous liquid, as opposed to the usual soft solid. If my memory doesn't fail me, my grandma used to keep her animals on a mostly starch-based diet, and still was able to process them for (solid) lard.

Can somebody confirm or disprove this theory?

Without going into the biochemical details of fatty acid digestion, mobilization, and synthesis, the brief answer is that the animal regulates its own fatty acid composition to within a specified range regardless of its food. Lipid metabolism does not differ from the metabolism of other compounds, in that it is a highly regulated, enzyme-mediated process, which ultimately is genetically determined.

It must be stressed that a heterotroph assimilates its food, not as complete macromolecules, but as their base components. For example, the starch-fed chickens you described would digest the starch down completely to glucose, which would then be converted in the liver to fatty acids. These fatty acids would be, of course, those that the chicken had the enzymes to synthesize, which are consequently mostly saturated fatty acids of a certain length that happen to solidify when cooled.


Why do sulfur compounds smell? Many of the compounds of sulfur have a strong odor.  Hydrogen sulfide from rotten eggs, the mercaptans of a skunk, the odor compounds in onions and garlic, the bitter taste of brassicas (cabbage, Brussels sprouts, collard greens, and similar plants) are all sulfur compounds.  What is it about sulfur that makes this so common?

Nice question! It would, obviously, be a long answer if I talk about all types of sulfur compounds here. So I will take up just 2 examples to explain this.

Thiols: when we talk about sulfur compounds with foul smell, thiols (mostly) come to the top in the list. Thiols (R-SH) are a class of compounds famous for their smell. Some examples of thiols include:

compounds in onion and garlic (see this answer for more details).
a mix of compounds responsible for the smell of skunk's spray (Anderson et al, 1975)
(R)/(S)-3-methyl-3-sulfanylhexan-1-ol found in human sweat
(Methylthio)methanethiol found in male mouse urine

Some biologically important examples would include coenzyme-A, glutathione and cysteine. However, not all thiols have fowl smell. For example, furan-2-ylmethanethiol provides the aroma of roasted coffee while grapefruit mercaptan gives grapefruit its characteristic scent. You can see a complete list here.

After seeing some examples, lets come to the main point i.e. why thiols have such foul smell. Thiols are decomposition products of proteins. When amino acids, like cysteine and methionine, are decomposed, thiols are one of the last things to form. This would explain why it becomes important to smell and avoid them (you wouldn't want to eat a dead and decaying mammoth, and a blind hunter might not be able to tell if he has found a decaying mammoth if he can't smell such compounds). As a side note, this also explains why organic acids have sour taste and smell. Compounds like glucuronic acid, citric acid, oxalic acid, lactic acid, butyric acid, acetic acid, etc. are all decomposition products (if you know how vinegar is made traditionally). This is why your body tries to get rid of them (I saw that in a video, if you eat things with foul smell and sour taste (it was a Japanese dish probably), you'll suffer nausea and vomiting since your body will do all it can to prevent its harmful effects). See this Wikipedia page for more info.

Hydrogen Sulfide: H2S, or hydrogen sulfide, is another compounds famous for its smell. It smells like rotten eggs, and indeed it is because H2S is what rotten eggs release. The basic reason why H2S smells so bad is again the same. Hydrogen sulfide is one of the final products of proteolysis during decomposition (see same Wikipedia article). But there is another cause why H2S has such foul smell (and why I didn't just stop at thiols): H2S is itself toxic.

H2S has many biological signalling functions, similar to NO and CO. These include:

H2S serves as endotheium-derived relaxing factor (EDRF) and endothelium-derived hyperpolarizing factor (EDHF). In short, it acts as smooth muscle relaxant and vasodilator. It also increases the response of NMDA receptor and facilitates long term potentiation in the brain.
it is converted to sulfite by thiosulfate reductase, and further to thiosulfate and sulfate by sulfite reductase, in the mitochondria. Sulfate is excreted in the urine.
it acts on ATP-sensitive potassium channels in smooth muscles, and does the blood vessel-relaxing work in smaller blood vessels.
it blunts, reverses and promotes healing of diverse inflammatory reactions. A full list is available on this Wikipedia page.
the most pronounced effect of H2S is much similar to that of CO. It also binds to the iron in the mitochondrial cytochrome c oxidase and other mitochondrial cytochrome enzymes, thus preventing cellular respiration and causing death.

This would explain why external source of H2S would be so harmful since, being biologically such an active molecule, its amount needs to be controlled strictly. Now, if you want to know why decaying bodies would be dangerous, this Wikipedia article would be a good starting point.

EDIT: The bitter taste of some vegetables is because of a type of compounds called glucosinolates found in some plants. These compounds dissociate upon eating to form isothiocyanate and are thought to be a part of the plants' defense system (see this answer for more details). Interestingly, a similar type of compound, cucurbitacin, found in Cucurbitacae plants, does not contain sulfur, yet it is responsible for the bitter taste of these plants.


What are these lentil like eggs? (Utah, United States) I've found these (what appear to be eggs) attached to an apple tree in central Utah USA.

The eggs are less than 1/2cm in length, have a waxy appearance and are attached to the branch at one end. Visually and tactilely they are quite similar to lentils.

I thought they were Box Elder bug eggs, but after a google image search I am convinced that is not what they are. What are they?

They look just like katydid eggs.

Image for comparison from the Missouri Botanical Garden:


How are viruses isolated? I read several different articles about virology, including the Wikipedia article on viruses and none of them explain how a virus is isolated. Some talk of injecting a virus into a chicken egg, but this does not isolate the virus, it just reproduces it.

If a researcher filters the particles in a serum to get those of a particular size, it proves nothing, because every creature has many thousands of different viruses in their blood, so any sample filtered by size will potentially have hundreds of different viruses in it, all the same size range. How can they possibly be separated?

Apologia

This answer was prepared in response to a previous question about Measles virus, but the question was removed for reasons I am not party to but which I can guess. As I spent some time on it I thought I’d post it in reply to this question. It is specific for measles virus, although there are obviously aspects that apply to other viruses, and deals with characterization as well as isolation. I am not a virologist (although I have done some biochemical work on pig herpes virus) so the it was is based solely on my reading of the papers cited and of Field’s Virology. I welcome suggestions for corrections or additions. It may seem to some that the question does not justify an answer of this detail. However it may also be of interest to young contemporary students and researchers who may have difficulty in envisaging the technological limitations of post-war experimental science.

Original Isolation and Characterization

The Edmonston strain of measles virus was isolated in 1954 by Enders and Peebles from the blood of a child with measles, David Edmonston. Before this scientists had found that an agent which was small enough to pass through filters that retained bacteria (i.e. possibly a virus) could cause measles when inoculated into monkeys and could be passaged in eggs containing developing chick embryos — the best method available until then. What Enders and Peebles managed to do was to maintain the virus in cells growing in tissue culture (a far more convenient host experimentally) — initially primary human kidney cells, and then monolayer cultures of chick embryo cells.

They used light microscopy to detect the changes in the cells that demonstrated that they were infected (the original paper contains many photographs of the cytopathogenic effects on infected cells). They could passage the virus by inoculating uninfected cells with infected cells, and they also isolated some virus from the cells by breaking open the cells and using the technique of centrifugation to separate the (larger) virus particles from the soluble cell contents and debris. (Subsequently much improved purification procedures were developed, resulting in preparations suitable for electron microscopy.)

One of the main techniques for characterizing viruses at this time was immunology. Thus, they found that the tissue culture fluid reacted with the antibodies in serum from measles patients by the complement fixation test, showing that it contained measles antigen. Previously they had showed that the cytopathogenicity of their preparations was neutralized by such antibody-containing serum.

Their paper ends with a summary of the evidence that what they isolated was the viral agent responsible for measles, and the further evidence needed for conclusive proof.

Subsequent characterization

The subsequent 60 years saw a steady improvement in analytical method, so that characterization did not rely solely on observation of cell pathology and immunology (the methodological repertoire of which also developed immeasurably). Electron microscopy in the 60s allowed cellular infection to be followed in greater detail and, after improvements in purification, the structure of the virus itself to be observed. The nucleic acid, protein and lipid envelope components could also be analysed with such purified virus.

The emerging techniques included gel electrophoresis of proteins (to determine their size and relative abundance, and to effect small-scale preparations). Larger scale purification of the virus proteins was possible with purified preparations, allowing  amino acid sequencing, especially of the most abundant proteins (e.g. the nucleocapsid protein, N),  and cDNA cloning and sequencing technology led to the sequencing of the whole Edmonston virus genome by Bellini et al. in 1985. This latter resulted in a much greater understanding of the gene expression of the virus, and provided tools such as specific oligonucleotide probes for various genes and antibody reagents to detect specific proteins.

Comment on Poster’s specific questions

The crux of the argument in the question would seem to be:

  “If a researcher filters the particles in a serum…any sample filtered by size will potentially have hundreds of different viruses in it, all the same size range…”

The above addresses this only implicitly. To be explicit:

Using either the serum from patients infected with a specific virus, virus passaged in eggs or viruses grown in cell culture one is massively increasing the proportion of that virus to any contaminants.
It is not easy to grow viruses in cell culture. Therefore when conditions are found for one virus they are unlikely to be optimal for other (low-level) contaminants. Thus the actual passage in culture removes viral contaminants. The purification problem becomes one of separation from cell debris.

Modern Methods of Purifying Measles Virus

Below is a summary flow chart taken from the article Manufacture of Measles Viruses by Langfield et al..

One of the key methods in this proces, Tangential Flow Filtration, is also known as Cross-Flow Filtration and is described in this Wikipedia article.


Suitable writings for self-study in general biology/chemistry I have recently commenced a graduate education in mathematics. Since I am currently attending extra classes in programming, which are not included in my education, my institution will not let me attend further courses in biology. Hereby, my question is simple:

What are suitable writings for self-study in general biology/chemistry? By general, I mean somewhat of an university-level introduction of mentioned fields.

My level of knowledge is slightly above "advanced" high school chemistry/biology.

(I apologize for my bad english, I'm Norwegian.)

My personal favorite intro biology book is Campbell Biology. The new edition is a little pricy, but you can probably find one of the older editions for pretty cheap--my intro class used 9th edition, and I don't think much has changed other than some of the taxon names.

I'm not sure what branch of Bio you're interested in, but Campbell covers all of introductory cell biology, organismal/evolutionary and ecology. The "concept check" questions at the end of each chapter are fantastic open-ended responses, and I believe most (if not all) have answers in the back of the book

I'm not as solid in general chemistry, but some good ones I remember are Brown, Lemay, and Burston which got me through AP Chem, and Zumdhal which my university uses and I haven't heard any complaints. The Brown text is a pretty good intro though. It provides good information and I thought it was pretty enjoyable to read.


Can we apply NGS (next generation sequencing) on DNA extracted from formalin fixed paraffin embedded tissues (FFPE)? I'm working on the microbiome and I can only use tissues which are fixed with formalin and embedded in paraffin. Is it possible to apply NGS (next generation sequencing) on FFPE (formalin fixed paraffin embedded) tissues? I'm particularly concerned about contamination that could affect the results.

It's a little hard to tell based on your question, but it seems that you want to do Illumina sequencing of FFPE tissue samples. This is quite routine, Illumina sells kits for it.
I do not know how the FFPE process affects microbial cells, but people have published recently on doing metagenomic analysis of FFPE tissue samples. This suggests to me that it is possible.
If "any contamination" is a problem, then maybe there are additional roadblocks; low-level contamination is a routine problem in metagenomics, which is exacerbated by complicated, old, or low-input sample types. I would suggest treating it the way you would treat any other metagenomic sample, and try to exclude the possibility of contamination through methods and account for it in data analysis.


Breathing's Role in Aerobic Respiration It's my understanding that aerobic cellular respiration is summarized in the formula $C_6H_{12}O_6+O_2→CO_2+H_2O+ATP$ (unbalanced). The glucose is thrown through a series of chemical reactions, ultimately resulting in the binding of a phosphate group to ADP to form ATP to store energy, and the carbon dioxide is a byproduct of this process. At the very end of this chain, though, there's a few electrons left hanging, and the oxygen is used to bind these to loose protons to form water.

It is also my understanding that the primary purposes of the lungs are to intake $O_2$ and exhale $CO_2$ for this process and their role in speech.

Assuming this is all correct, are the secondary, minor functions of the lungs important enough that if there was a way to siphon the $e^-$ and $CO_2$ to some other function and find something else to allow the larynx to work we wouldn't need to breathe?

There's a fairly large misconception here. The huge amount of energy that's available from respiration, and which drives the synthesis of ATP, in fact comes from that reaction that you consider secondary, and that's the transfer of electrons to oxygen.

If you want to think of it as two separate reactions, then the first one
$$Glucose + 12\ {NAD}^+ + 6 H2O ⇌ 6 CO_2 + 12 NADH$$ releases about 10 times less energy than the second one: $$6 O_2 + 12 NADH ⇌ 12 NAD^+ + 12 H_2O$$

The reason you get something resembling 30 ATP from a glucose molecule is because oxygen is such a great electron acceptor.

Many microbes can in fact use molecules other than oxygen as the electron acceptor, but in general it is much less energetically efficient than using oxygen.

You could try equilibrator for estimating the free energy available from various chemical reactions. First reaction and second reaction from above as examples.


How does 'phosphorylation of glucose' maintain concentration gradient in membrane transport (facilitated diffusion)? Our book explains how glucose from the blood plasma gets inside red blood cells via facilitated transport.

It states here in the book that the glucose will be transported inside by a carrier protein. Then the glucose will be phosphorylated by ATP. The phosphorylation of glucose molecules maintains the concentration gradient. This prevents the glucose from diffusing back to the blood plasma.

Based from my current understanding, concentration only change when the number of dissolved molecules changes.

My question is, "How can phosphorylation maintain the concentration gradient if it doesn't change the number of molecules inside the red blood cell?"

Sort of a "magic trick" that biology does.

With facilitated transport, the movement is passive. That is, you have a protein that certain molecules/ions can pass through, but you aren't doing any pumping or using energy to move molecules. That means that the molecules are free to move in either direction, according to concentration (and electrical) gradients.

Therefore, assuming no influence of electrical gradients, the highest concentration a passively transported molecule can reach is equal concentration inside and out.

Chemically, phosphorylated glucose is not glucose. It is phosphorylated glucose. If you phosphorylate every glucose molecule that comes in, your internal concentration of glucose would stay at zero. If phosphorylated glucose can't leave via the same facilitated diffusion pathway, then glucose will effectively only come in via that pathway, it can't go back out. The result is that, although glucose inside can never be higher than glucose outside, the sum total of phosphorylated glucose + glucose inside can continue to rise above the concentration of glucose outside.

Of course, this step isn't exactly free, because it costs the energy of phosphorylation.


Can a constitutively active kinase be highly regulated? I am studying the protein kinase GSK3 and I am learning about the regulation of its activity. Many journal papers that I have read have stated that GSK3 is unique because it is a constitutively active kinase, meaning that it is always active and is not activated in response to a specific stimulus/condition. I have read that phosphorylation of GSK3 at specific tyrosine residues only enhances its activity, but the GSK3 itself is always active.
However, many of these papers also talk about the various mechanisms through which the activity of GSK3 is regulated (e.g. through phosphorylation at specific amino acid residues, subcellular localisation, protein complex formation). One paper states the following:

The many functions regulated by GSK3 (through its phosphorylation of
numerous substrates) suggests that the activity of GSK3 must be highly
regulated. Four key mechanisms have been identified that contribute to
regulating the actions of GSK3 in a substrate-specific manner. These include regulation by phosphorylation of GSK3 itself, the subcellular localization of GSK3, the formation of protein complexes containing GSK3, and the phosphorylation state of GSK3 substrates.

I am wondering how is it possible for a kinase to be constitutively active and be highly regulated at the same time? If it is constitutively active, doesn't this imply that it is always active? Any insights are appreciated.

I think you have answered your own question, but yes, the nomenclature is confusing. When they say it is a constitutively active kinase, you are right, it plainly means the enzyme can always phosphorylate other molecules/proteins.
The regulation part is what you have in the review paper you cited. If the kinase is always on but localized in a compartment where there is nothing to phosphorylate, it could be considered "regulated". The same logic applies to the other regulation methods listed.


Why are vegetations in Infective endocarditis common on the atrial side? Robbin's Pathology says that vegetations of IE are more common on the atrial side in AV valves. In Liebmann Sack's Endocarditis, which is a sterile (non bacterial) type of endocarditis, the underside of the valve also has vegetations. Why is it so? If its just colonization of bacteria in IE, how does it matter where?

Does gravity matter in colonisation of bacteria!! ?

With the bacterial IE, it would look like the vegetations tend to colonize the lower pressure sides of the valve structure (REF1). Something that was made clear that distinguishes the IE vegetations frm the Libman-Sacks vegetations was overall motility. Specifically, Libman-Sacks vegetations exhibit no independent motion and are largely sessile, whereas IE vegetations can be motile (REF2). In terms of where they attach I think this makes sense, the more motile vegetations colonize the less turbulent area, and the sessile vegetations seem to hold up in either scenario. Reference 1 also rather usefully points out that IE vegetations may be located anywhere on the valve structure as well, they just tend to be on the low-pressure side.

So being said i would argue that the motility of the organisms is a deciding factor.


Can bacteria be killed by purely physical trauma? In a question over on Skeptics Stack Exchange, I half-jokingly listed "blunt force trauma" as a means by which a bacterium could be killed. That makes me wonder if it really is only a joke or whether such an assault is plausible. The webcomic xkcd also did an episode on the possibility of physically attacking a cell (though not a bacterial one).
There are plenty of resources on how bacteria can die from biological processes (e.g. infection with a bacteriophage, biological effects of antibiotics, etc.), chemical attacks (e.g. membranes getting lysed with bleach), temperature (i.e. cooking), and radiation, but have bacteria ever been shown to be vulnerable to a direct physical assault such as being hit with a hammer, cut with a knife, thrown onto a hard surface, twisted until they burst, or being drawn and quartered?
I do know that many common weapons that are generally effective against animal life (swords, axes, clubs, assault rifles, etc.) do not have the exact same properties at the microscopic scale, but does that mean that physically beating a bacterium to death is absurd or is there a known way to kill a bacterium using only physical force?
To be clear, I recognize that our normal means of killing unwanted bacteria - cooking food to prevent food poisoning, treating bacterial infections with antibiotics, cleaning hard surfaces with bleach, applying rubbing alcohol to wounds, etc., are probably effective enough that most people don't bother pulling out that old Wild West six-shooter they have in the closet or the antique sword hanging on their wall the way they might consider them if attacked by a bear or wolf, but that's not the question.

My (limited) understanding is that it is quite hard to avoid killing some bacteria even with very gentle physical manipulation. On the other hand, it is quite hard to use physical force to achieve reasonable level of sterilization. Let's bring some examples with a few (hastily found) references.
Pressure
My guess is that most examples the OP mentions (hit with a hammer, knife) will exert their biggest effect on bacteria through a peak in pressure - bacteria are comparable in size to width of a moderately sharp knife edge and thus there will be more "crushing" than "cutting" even with a knife (Science of Sharp discusses that Chosera 1k - a very good stone - will sharpen a knife to 0.5 μm, Wiki has bacteria at 0.5 - 5 μm).
I recently collaborated on a project where initial interpretation of results was wrong simply because putting a glass slip over a sample to inspect it under a microscope caused around 10% of bacterial (B. subtilis) cells to die. Putting a small weight on the slip resulted in roughly 80kPa additional pressure applied and ~ 80% of the cells dying within 25 minutes. But 80% death rate won't do much to e.g. reduce risk of food poisoning - as long as there are nutrients, bacteria will regrow quickly.
As m4rio mentions in a comment to another answer, and in line with The effect of high pressure processing on the microbial, physical and chemical properties of Valencia and Navel orange juice you need to ramp up the pressure to 100s of MPa to achieve high efficacy. And even that may not work so cleanly - e.g., the linked paper mentions that after applying 600 MPa, bacteria regrew to detectable levels after 4 weeks of storage at 4°C. Killing all (or almost all) bacteria is just really hard.
Because it was fun, I'll also link to study showing that High pressure-processed guacamole (4 cycles, 689 MPa, 5 min each) is clean enough to not spoil for 30 days in 25°C.
Note: the two sources above were linked from a commercial vendor of high-pressure products and may thus be biased. While the exact numbers might be problematic, I think the main point is unlikely to be affected.
A review High-pressure processing – effects on microbial food safety and food quality mentions some limitations, e.g.:

Endospores tend to be extremely HPP resistant compared with vegetative cells, withstanding treatments of more than 1000 MPa

ascospores of heat-resistant moulds such as Byssochlamys, Neosartorya and Talaromyces are generally considered to be extremely HPP resistant

It should however be noted that it seems that high-pressure treatment is good enough to be semi-commonly used as the main sterilization for suitable foodstuffs - Current status and future trends of high-pressure processing in food industry reports high-pressure treatment to be FDA approved and USD 10 billion worth of food being treated as of 2015.
Pressure for killing cells previously briefly discussed at Can bacteria or other microorganisms be killed by applying pressure?
Ultrasound
While frequently used for cell lysis in lab, it also doesn't kill all cells. For example The Sterilization of Suspensions Contaminated with MicroorganismsUsing Ultrasound Irradiation reports that roughly 20-30% of bacteria (Bacillus) survived 60 minutes of ultrasound treatment.
Ball mill
This is probably closest to actual "physical trauma" - you put small beads and a suspension with microorganisms into a cylinder and roll it hard. A low energy process for the recovery of bioproducts from cyanobacteria using a ball mill discusses using this process to extract proteins from cyanobacteria. They don't report absolute number of cells killed, just that it varies considerably with the rotation speed, size of beads, .... My interpretation is that in most cases a non-negligible proportion of the cells survive while majority is killed (and their contents extracted for further use).


What does "melanistic phase" mean in regards to black jaguars? From Black Jaguars in Belize?: A survey of melanism in the jaguar, panthera onca by John R. Meyer:

  There occurs in jaguars, as well as in leopards (Panthera pardus), a melanistic phase which has been much sought by zoos for exhibition purposes. This color phase was relatively difficult to acquire until the 1970's, by which time a few zoos, among them the Jacksonville Zoo, had begun to breed them regularly.

And likewise on the Wikipedia page for Black Panther:

  This "black cougar" was most likely a margay or ocelot, which are under 18 kg (40 lb) in weight, live in trees, and do have melanistic phases.

I don't understand this term "melanistic phase", which makes it sound like the jaguar has a period of time in which it is black.  However, it looks like it means something else, i.e., along the lines of What does phasing mean?, but I don't really understand that.

What does "melanistic phase" mean in regards to black jaguars?

In this context, "phase" means the same thing as "color phase". Here's the definition from Merriam-Webster:

  1: a seasonally variant pelage color
  2: (a) a genetic variant manifested by the occurrence of a skin or pelage color unlike the wild type of the animal group in which it appears (b) an individual marked by such a variant

So it sounds as if "phase" originated in the case of animals that spent a season with a particular pelage (meaning "fur, hair or wool"), such as the snowshoe hare, which is brown in summer and turns white in winter, and the term then began to be used for animals that might be a different color from other individuals of that species. If the alternate color was caused by a recessive allele, most offspring of that individual would have the wildtype color unless bred with another double-recessive individual, and alternate offspring might spontaneously be born to wildtype heterozygous individuals, so pre-Mendel biologists might have seen it as a "phase" between wildtype generations.

As you can see from the illustrations on the snowshoe hare Wikipedia article, "morph" is much more commonly used today for this kind of visible variation within a species. Wikipedia has a fuller description of the nomenclature of this and related terms.


Can citric acid be used as fungicide in banana based fly (Drosophila melanogaster) food? I am experimenting with a banana based fly food for better larval synchronization. Instead of propionic acid, I am going to use (added) citric acid as fungicide. Is there any downside to this approach?

I have no data on the fungicidal and insecticidal properties of citric acid at various concentrations (and you may be hard pressed to find a good study on this), but I have worked in some fly labs and I understand the composition of fly food only peripherally. The most commonly and most abundantly used fungicide in a lot of fly labs is methylparaben (trade names Tegosept or Nipagin in Europe). I imagine food recipe-makers have arrived at this solution through trial and error over the years; this is as much an art as it is a science, in balancing the ingredients with humidity, temperature, species and genotype with desired generation time. Methylparaben is used with and without propionic acid for Drosophila melanogaster. Do you have methylparaben in your fly food? Also, I think citric acid is not commonly used, again I imagine because technicians probably found it inferior to its alternatives.
The National Drosophila Species Stock Center at Cornell uses methylparaben with its banana food recipe.
Alternatively, please take a look at the variety of recipes Bloomington Drosophila Stock Center uses. To my knowledge they supply and stock the largest collection using the most recipes.

Fly foods
Alternative recipes


Can toilet or sewage aerosol transfer eggs of human parasites? Flushing a toilet produces hundreds of thousands of tiny droplets containing viruses and bacteria (source).
Additionally vomit and feces of infected people sometimes contain extremely high concentrations of viruses or bacteria.
Similarly, sewage water produces droplets and was proposed as the transmission mode in the SARS outbreak in Hong Kong in 2003.

Apart from viruses and bacteria being transferred through toilet and sewage aerosol, can human parasite eggs be transferred as well given the droplet size?
I'm interested in the aerosol, not the relatively large droplets that land on the toilet itself.

Common small protozoa that are present in infective loads in stool have not been documented in toilet aerosols. The size of the organism is going to be an issue here. As you can see from this table from the CDC, intestinal parasites (and ova) of concern to human health are larger than most droplets. Compare, E. coli and Shigella (0.5 - ~2 $\mu m$). Salmonella tends to be a litter larger (up to 5$\mu m$), but I expect those 5 $\mu m$ specimens are not the ones transmitted through aerosols. Cryptosporidium would probably be the most likely candidate, and that's one shown in the first linked study to not be present in aerosols. I wouldn't say it's impossible, but certainly transmission through aerosolized particles is less likely than for the smaller viral and bacterial organisms.


What is the mechanism responsible for the periodicity of hiccups? Hiccups (singultus) are caused by involuntary spasm of the intercostal muscles and diaphragm, producing a sudden inspiration with associated glottic closure. It seems the neurobiology of the reflex has been relatively well worked out (nicely summarized in this review):

Afferent limb: the vagus, phrenic, and cervical sympathetic chain
Central mediators: the hypothalamus, medullary reticular formation, and brainstem respiratory centers
Efferent limb:  primarily mediated by the phrenic nerve but also including the intercostals, and secondarily the recurrent laryngeal nerve to cause glottic closure

To my knowledge, hiccups are always cyclical, with a period of between 1 and 30 seconds. According to the afore cited article, this interval is relatively constant it a given individual. The cycling stops after some time; efferent vagal tone can suppress it, or it stops spontaneously. My question is:

What is the mechanism responsible for this periodicity?

Something must be triggering that afferent limb every 1 - 30s, but I can’t think of a normal physiologic process that triggers the firing of one or more nerves with such periodicity and then stops abruptly.

Check out Van Der Pol (chaotic) oscillator. It can be used to model heart rhythms. If a stimulus is applied to the oscillator, it will return to the curved envelope discussed in the article.

However, if "pushed just right", the oscillator's current state will be pushed towards the very center of the loop, and it will stop oscillating. This is the mechanism by which a sudden blow to the chest can stop the heart.

Reading up on hiccups, their "almost regular" interval sounds similar to a chaotic oscillator. The abrupt stop may represent pushing the oscillator towards the middle of the limit cycle, from which it cannot recover, effectively stopping hiccups


Which organism built this 2 cm long sandy tube? I found this tube and wondered which animal built it.

It is 2 cm long and has an entrance/exit hole on one side. The material seems to be sandy.

I assume it is for breeding purposes.

Before finding the tube, I saw an unfamiliar kind of wasp around (could also be a hornet), with very long black legs. However, I could not find a wasp or hornet that builds such tubes, so it could be unrelated.

Location is Germany.

My questions are: which organism built this structure and should I expect anything crawling out one day? Or is it unfinished, since the entrance isn't closed?

This is the nest of a Mud dauber, also known as Mud wasp. This was possibly made by a Black and Yellow Mud dauber based on the following information.

The nest of the black and yellow mud dauber is a simple, one-cell, urn-shaped nest. 1*

As for expecting something coming out of it? I doubt it, as the nests are sealed after depositing an egg.

After building a cell of the nest, the female wasp captures several spiders. The captured prey are stung and paralyzed before being placed in the nest (usually 6-15 per cell), and then a single egg is deposited on the prey within each cell. The wasp then seals the cell with a thick mud plug. 2*
Eventually, the hatching larva will eat the prey and emerge from the nest.