Text-to-Image
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- ---
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- license: apache-2.0
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- ---
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
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+ ---
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+ license: apache-2.0
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+ ---
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+ language:
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+ - en
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+ pipeline_tag: text-to-image
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+ ---
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+ # BizGen: Advancing Article-level Visual Text Rendering for Infographics Generation
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+
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+ <a href="https://arxiv.org/abs/2406.04314"><img src="https://img.shields.io/badge/Paper-arXiv-red?style=for-the-badge" height=22.5></a>
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+ <a href="https://github.com/RockeyCoss/SPO"><img src="https://img.shields.io/badge/Gihub-Code-succees?style=for-the-badge&logo=GitHub" height=22.5></a>
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+ <a href="https://rockeycoss.github.io/spo.github.io/"><img src="https://img.shields.io/badge/Project-Page-blue?style=for-the-badge" height=22.5></a>
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+
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+ <table>
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+ <tr>
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+ <td><img src="assets/imgs/0.png" alt="teaser example 0" width="200"/></td>
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+ <td><img src="assets/imgs/1.png" alt="teaser example 1" width="200"/></td>
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+ <td><img src="assets/imgs/2.png" alt="teaser example 2" width="200"/></td>
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+ <td><img src="assets/imgs/3.png" alt="teaser example 3" width="200"/></td>
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+ </tr>
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+ </table>
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+
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+ ## Abstract
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+ <p>
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+ Generating visually appealing images is fundamental to modern text-to-image generation models.
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+ A potential solution to better aesthetics is direct preference optimization (DPO),
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+ which has been applied to diffusion models to improve general image quality including prompt alignment and aesthetics.
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+ Popular DPO methods propagate preference labels from clean image pairs to all the intermediate steps along the two generation trajectories.
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+ However, preference labels provided in existing datasets are blended with layout and aesthetic opinions, which would disagree with aesthetic preference.
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+ Even if aesthetic labels were provided (at substantial cost), it would be hard for the two-trajectory methods to capture nuanced visual differences at different steps.
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+ </p>
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+ <p>
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+ To improve aesthetics economically, this paper uses existing generic preference data and introduces step-by-step preference optimization
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+ (SPO) that discards the propagation strategy and allows fine-grained image details to be assessed. Specifically,
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+ at each denoising step, we 1) sample a pool of candidates by denoising from a shared noise latent,
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+ 2) use a step-aware preference model to find a suitable win-lose pair to supervise the diffusion model, and
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+ 3) randomly select one from the pool to initialize the next denoising step.
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+ This strategy ensures that diffusion models focus on the subtle, fine-grained visual differences
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+ instead of layout aspect. We find that aesthetic can be significantly enhanced by accumulating these
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+ improved minor differences.
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+ </p>
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+ <p>
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+ When fine-tuning Stable Diffusion v1.5 and SDXL, SPO yields significant
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+ improvements in aesthetics compared with existing DPO methods while not sacrificing image-text alignment
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+ compared with vanilla models. Moreover, SPO converges much faster than DPO methods due to the step-by-step
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+ alignment of fine-grained visual details.
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+ </p>
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+
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+ ## Model Description
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+
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+ This model is fine-tuned from [stable-diffusion-xl-base-1.0](https://huggingface.co/stabilityai/stable-diffusion-xl-base-1.0). It has been trained on 4,000 prompts for 10 epochs.
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+
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+ This is a merged checkpoint that combines the LoRA checkpoint with the base model [stable-diffusion-xl-base-1.0](https://huggingface.co/stabilityai/stable-diffusion-xl-base-1.0). If you want to access the LoRA checkpoint, please visit [SPO-SDXL_4k-p_10ep_LoRA](https://huggingface.co/SPO-Diffusion-Models/SPO-SDXL_4k-p_10ep_LoRA). We also provide a LoRA checkpoint compatible with [stable-diffusion-webui](https://github.com/AUTOMATIC1111/stable-diffusion-webui), which can be accessed [here](https://civitai.com/models/510261?modelVersionId=567119).
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+
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+
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+ ## A quick example
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+ ```python
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+ from diffusers import StableDiffusionXLPipeline, UNet2DConditionModel
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+ import torch
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+
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+ # load pipeline
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+ inference_dtype = torch.float16
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+ pipe = StableDiffusionXLPipeline.from_pretrained(
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+ "SPO-Diffusion-Models/SPO-SDXL_4k-p_10ep",
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+ torch_dtype=inference_dtype,
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+ )
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+ vae = AutoencoderKL.from_pretrained(
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+ 'madebyollin/sdxl-vae-fp16-fix',
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+ torch_dtype=inference_dtype,
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+ )
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+ pipe.vae = vae
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+ pipe.to('cuda')
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+
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+ generator=torch.Generator(device='cuda').manual_seed(42)
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+ image = pipe(
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+ prompt='a child and a penguin sitting in front of the moon',
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+ guidance_scale=5.0,
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+ generator=generator,
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+ output_type='pil',
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+ ).images[0]
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+ image.save('moon.png')
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+ ```
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+
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+ ## Citation
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+ If you find our work or codebase useful, please consider giving us a star and citing our work.
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+ ```
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+ @article{liang2024step,
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+ title={Aesthetic Post-Training Diffusion Models from Generic Preferences with Step-by-step Preference Optimization},
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+ author={Liang, Zhanhao and Yuan, Yuhui and Gu, Shuyang and Chen, Bohan and Hang, Tiankai and Cheng, Mingxi and Li, Ji and Zheng, Liang},
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+ journal={arXiv preprint arXiv:2406.04314},
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+ year={2024}
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+ }
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+ ```