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README.md
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---
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<style>
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.title {
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font-size:
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text-align: center;
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color: #333;
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font-family: Arial, sans-serif;
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text-transform: uppercase;
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letter-spacing: 0.05em;
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padding: 0.5em 0;
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background: transparent;
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box-shadow: 0px 0px 20px 0px rgba(0,0,0,0.15);
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display: inline-block;
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width: auto;
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}
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.title span {
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background: -webkit-linear-gradient(45deg, #fe6b8b 30%, #ff8e53 90%);
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</div>
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</div>
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## Overview
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Hermitage XL is a high-resolution, latent text-to-image diffusion model. The model has been fine-tuned using a learning rate of 4e-7 over 5000 steps with a batch size of 16 on a curated dataset of superior-quality anime-style images. This model is derived from Stable Diffusion XL 1.0.
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- Use it with 🧨 [`diffusers`](https://huggingface.co/docs/diffusers/index)
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- Use it with the [`ComfyUI`](https://github.com/comfyanonymous/ComfyUI)
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## Features
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1. High-Resolution Images: The model trained with 1024x1024 resolution. The model is trained using [NovelAI Aspect Ratio Bucketing Tool](https://github.com/NovelAI/novelai-aspect-ratio-bucketing) so that it can be trained at non-square resolutions.
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2. Anime-styled Generation: Based on given text prompts, the model can create high quality anime-styled images.
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3. Fine-Tuned Diffusion Process: The model utilizes a fine-tuned diffusion process to ensure high quality and unique image output.
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## Model Details
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- **Developed by:** [Linaqruf](https://github.com/Linaqruf)
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- **Model Description:** This is a model that can be used to generate and modify anime-themed images based on text prompts.
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- **License:** [CreativeML Open RAIL++-M License](https://huggingface.co/stabilityai/stable-diffusion-2/blob/main/LICENSE-MODEL)
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- **Finetuned from model:** [Stable Diffusion XL 1.0](https://huggingface.co/stabilityai/stable-diffusion-xl-base-1.0)
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## How to Use:
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- Download `Hermitage XL` [here](https://huggingface.co/Linaqruf/hermitage-xl/resolve/main/hermitage-xl.safetensors), the model is in `.safetensors` format.
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```
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masterpiece, best quality, illustration, beautiful detailed, finely detailed, dramatic light, intricate details
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```
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## 🧨 Diffusers
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image.save("anime_girl.png")
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```
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## Limitation
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1. This model inherit Stable Diffusion XL 1.0 [limitation](https://huggingface.co/stabilityai/stable-diffusion-xl-base-1.0#limitations)
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2. This model is overfitted and cannot follow prompts well, because it's fine-tuned for 5000 steps with small scale datasets.
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3. It's only a preview model to find good hyperparameter and training config for Stable Diffusion XL 1.0
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## Example
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Here is some cherrypicked samples and comparison between available models:
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</td>
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</tr>
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</table>
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<hr>
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---
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<style>
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.title-container {
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display: flex;
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justify-content: center;
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align-items: center;
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height: 100vh; /* Adjust this value to position the title vertically */
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}
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.title {
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font-size: 3em;
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text-align: center;
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color: #333;
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font-family: Arial, sans-serif;
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text-transform: uppercase;
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letter-spacing: 0.05em;
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padding: 0.5em 0;
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box-shadow: 0px 0px 20px 0px rgba(0,0,0,0.15);
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background: transparent;
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}
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.title span {
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background: -webkit-linear-gradient(45deg, #fe6b8b 30%, #ff8e53 90%);
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</div>
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</div>
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<hr>
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## Overview
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Hermitage XL is a high-resolution, latent text-to-image diffusion model. The model has been fine-tuned using a learning rate of 4e-7 over 5000 steps with a batch size of 16 on a curated dataset of superior-quality anime-style images. This model is derived from Stable Diffusion XL 1.0.
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- Use it with 🧨 [`diffusers`](https://huggingface.co/docs/diffusers/index)
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- Use it with the [`ComfyUI`](https://github.com/comfyanonymous/ComfyUI)
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<hr>
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## Features
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1. High-Resolution Images: The model trained with 1024x1024 resolution. The model is trained using [NovelAI Aspect Ratio Bucketing Tool](https://github.com/NovelAI/novelai-aspect-ratio-bucketing) so that it can be trained at non-square resolutions.
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2. Anime-styled Generation: Based on given text prompts, the model can create high quality anime-styled images.
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3. Fine-Tuned Diffusion Process: The model utilizes a fine-tuned diffusion process to ensure high quality and unique image output.
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<hr>
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## Model Details
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- **Developed by:** [Linaqruf](https://github.com/Linaqruf)
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- **Model Description:** This is a model that can be used to generate and modify anime-themed images based on text prompts.
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- **License:** [CreativeML Open RAIL++-M License](https://huggingface.co/stabilityai/stable-diffusion-2/blob/main/LICENSE-MODEL)
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- **Finetuned from model:** [Stable Diffusion XL 1.0](https://huggingface.co/stabilityai/stable-diffusion-xl-base-1.0)
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<hr>
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## How to Use:
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- Download `Hermitage XL` [here](https://huggingface.co/Linaqruf/hermitage-xl/resolve/main/hermitage-xl.safetensors), the model is in `.safetensors` format.
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```
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masterpiece, best quality, illustration, beautiful detailed, finely detailed, dramatic light, intricate details
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```
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<hr>
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## 🧨 Diffusers
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image.save("anime_girl.png")
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```
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<hr>
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## Limitation
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1. This model inherit Stable Diffusion XL 1.0 [limitation](https://huggingface.co/stabilityai/stable-diffusion-xl-base-1.0#limitations)
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2. This model is overfitted and cannot follow prompts well, because it's fine-tuned for 5000 steps with small scale datasets.
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3. It's only a preview model to find good hyperparameter and training config for Stable Diffusion XL 1.0
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<hr>
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## Example
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Here is some cherrypicked samples and comparison between available models:
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</td>
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</tr>
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</table>
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